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868953 | https://en.wikipedia.org/wiki/Congenital%20heart%20defect | Congenital heart defect | A congenital heart defect (CHD), also known as a congenital heart anomaly, congenital cardiovascular malformation, and congenital heart disease, is a defect in the structure of the heart or great vessels that is present at birth. A congenital heart defect is classed as a cardiovascular disease. Signs and symptoms depend on the specific type of defect. Symptoms can vary from none to life-threatening. When present, symptoms are variable and may include rapid breathing, bluish skin (cyanosis), poor weight gain, and feeling tired. CHD does not cause chest pain. Most congenital heart defects are not associated with other diseases. A complication of CHD is heart failure.
Congenital heart defects are the most common birth defect. In 2015, they were present in 48.9 million people globally. They affect between 4 and 75 per 1,000 live births, depending upon how they are diagnosed. In about 6 to 19 per 1,000 they cause a moderate to severe degree of problems. Congenital heart defects are the leading cause of birth defect-related deaths: in 2015, they resulted in 303,300 deaths, down from 366,000 deaths in 1990.
The cause of a congenital heart defect is often unknown. Risk factors include certain infections during pregnancy such as rubella, use of certain medications or drugs such as alcohol or tobacco, parents being closely related, or poor nutritional status or obesity in the mother. Having a parent with a congenital heart defect is also a risk factor. A number of genetic conditions are associated with heart defects, including Down syndrome, Turner syndrome, and Marfan syndrome. Congenital heart defects are divided into two main groups: cyanotic heart defects and non-cyanotic heart defects, depending on whether the child has the potential to turn bluish in color. The defects may involve the interior walls of the heart, the heart valves, or the large blood vessels that lead to and from the heart.
Congenital heart defects are partly preventable through rubella vaccination, the adding of iodine to salt, and the adding of folic acid to certain food products. Some defects do not need treatment. Others may be effectively treated with catheter based procedures or heart surgery. Occasionally a number of operations may be needed, or a heart transplant may be required. With appropriate treatment, outcomes are generally good, even with complex problems.
Signs and symptoms
Signs and symptoms are related to type and severity of the heart defect. Symptoms frequently present early in life, but it is possible for some CHDs to go undetected throughout life. Some children have no signs while others may exhibit shortness of breath, cyanosis, fainting, heart murmur, under-development of limbs and muscles, poor feeding or growth, or respiratory infections. Congenital heart defects cause abnormal heart structure resulting in production of certain sounds called heart murmur. These can sometimes be detected by auscultation; however, not all heart murmurs are caused by congenital heart defects.
Associated conditions
Congenital heart defects are associated with an increased incidence of seven other specific medical conditions, together being called the VACTERL association:
V — Vertebral anomalies
A — Anal atresia
C — Cardiovascular anomalies
T — Tracheoesophageal fistula
E — Esophageal atresia
R — Renal (Kidney) and/or radial anomalies
L — Limb defects
Ventricular septal defect (VSD), atrial septal defect (ASD), and tetralogy of Fallot (ToF) are the most common congenital heart defects seen in the VACTERL association. Less common defects in the association are persistent truncus arteriosus and transposition of the great arteries.
Causes
The cause of congenital heart disease may be genetic, environmental, or a combination of both.
Genetic
Genetic mutations, often sporadic, represent the largest known cause of congenital heart defects. They are described in the table below.
Molecular pathways
The genes regulating the complex developmental sequence have only been partly elucidated. Some genes are associated with specific defects. A number of genes have been associated with cardiac manifestations. Mutations of a heart muscle protein, α-myosin heavy chain (MYH6) are associated with atrial septal defects. Several proteins that interact with MYH6 are also associated with cardiac defects. The transcription factor GATA4 forms a complex with the TBX5 which interacts with MYH6. Another factor, the homeobox (developmental) gene, NKX2-5 also interacts with MYH6. Mutations of all these proteins are associated with both atrial and ventricular septal defects; In addition, NKX2-5 is associated with defects in the electrical conduction of the heart and TBX5 is related to the Holt–Oram syndrome which includes electrical conduction defects and abnormalities of the upper limb. The Wnt signaling co-factors BCL9, BCL9L and PYGO might be part of these molecular pathways, as when their genes are mutated, this causes phenotypes similar to the features present in Holt-Oram syndrome. Another T-box gene, TBX1, is involved in velo-cardio-facial syndrome DiGeorge syndrome, the most common deletion which has extensive symptoms including defects of the cardiac outflow tract including tetralogy of Fallot.
The notch signaling pathway, a regulatory mechanism for cell growth and differentiation, plays broad roles in several aspects of cardiac development. Notch elements are involved in determination of the right and left sides of the body plan, so the directional folding of the heart tube can be impacted. Notch signaling is involved early in the formation of the endocardial cushions and continues to be active as the develop into the septa and valves. It is also involved in the development of the ventricular wall and the connection of the outflow tract to the great vessels. Mutations in the gene for one of the notch ligands, Jagged1, are identified in the majority of examined cases of arteriohepatic dysplasia (Alagille syndrome), characterized by defects of the great vessels (pulmonary artery stenosis), heart (tetralogy of Fallot in 13% of cases), liver, eyes, face, and bones. Though less than 1% of all cases, where no defects are found in the Jagged1 gene, defects are found in Notch2 gene. In 10% of cases, no mutation is found in either gene. For another member of the gene family, mutations in the Notch1 gene are associated with bicuspid aortic valve, a valve with two leaflets instead of three. Notch1 is also associated with calcification of the aortic valve, the third most common cause of heart disease in adults.
Mutations of a cell regulatory mechanism, the Ras/MAPK pathway are responsible for a variety of syndromes, including Noonan syndrome, LEOPARD syndrome, Costello syndrome and cardiofaciocutaneous syndrome in which there is cardiac involvement. While the conditions listed are known genetic causes, there are likely many other genes which are more subtle. It is known that the risk for congenital heart defects is higher when there is a close relative with one.
Environmental
Known environmental factors include certain infections during pregnancy such as rubella, drugs (alcohol, hydantoin, lithium and thalidomide) and maternal illness (diabetes mellitus, phenylketonuria, and systemic lupus erythematosus). Alcohol exposure in the father also appears to increase the risk of congenital heart defects.
Being overweight or obese increases the risk of congenital heart disease. Additionally, as maternal obesity increases, the risk of heart defects also increases. A distinct physiological mechanism has not been identified to explain the link between maternal obesity and CHD, but both pre-pregnancy folate deficiency and diabetes have been implicated in some studies.
Twins and Multiple Births
Congenital heart defects happen more often in twins than in single babies. Monochorionic twins, who share a placenta, have a greater risk of these heart defects compared to dichorionic twins, who have their own placentas. A systematic review and meta-analysis of four studies conducted in 2007 showed a 9-fold increase in CHD risk in MC twins compared to singletons.
Mechanism
There is a complex sequence of events that result in a well formed heart at birth and disruption of any portion may result in a defect. The orderly timing of cell growth, cell migration, and programmed cell death ("apoptosis") has been studied extensively and the genes that control the process are being elucidated.
Around day 15 of development, the cells that will become the heart exist in two horseshoe shaped bands of the middle tissue layer (mesoderm), and some cells migrate from a portion of the outer layer (ectoderm), the neural crest, which is the source of a variety of cells found throughout the body. On day 19 of development, a pair of vascular elements, the "endocardial tubes", form. The tubes fuse when cells between then undergo programmed death and cells from the first heart field migrate to the tube, and form a ring of heart cells (myocytes) around it by day 21. On day 22, the heart begins to beat and by day 24, blood is circulating.
At day 22, the circulatory system is bilaterally symmetrical with paired vessels on each side and the heart consisting of a simple tube located in the midline of the body layout. The portions that will become the atria and will be located closest to the head are the most distant from the head. From days 23 through 28, the heart tube folds and twists, with the future ventricles moving left of center (the ultimate location of the heart) and the atria moving towards the head.
On day 28, areas of tissue in the heart tube begin to expand inwards; after about two weeks, these expansions (the membranous "septum primum" and the muscular "endocardial cushions") fuse to form the four chambers of the heart. A failure to fuse properly will result in a defect that may allow blood to leak between chambers. After this happens, cells that have migrated from the neural crest begin to divide the bulbus cordis. The main outflow tract is divided in two by the growth of a spiraling septum, becoming the great vessels—the ascending segment of the aorta and the pulmonary trunk. If the separation is incomplete, the result is a "persistent truncus arteriosus". The vessels may be reversed ("transposition of the great vessels"). The two halves of the split tract must migrate into the correct positions over the appropriate ventricles. A failure may result in some blood flowing into the wrong vessel (e.g. overriding aorta). The four-chambered heart and the great vessels have features required for fetal growth. The lungs are unexpanded and cannot accommodate the full circulatory volume. Two structures exist to shunt blood flow away from the lungs to compensate. Cells in part of the septum primum die, creating a hole while new muscle cells (the "septum secundum") grow along the right atrial side of the septum primum except for one region, leaving a gap through which blood can pass from the right atrium to the left atrium (the foramen ovale). A small vessel called the ductus arteriosus allows blood from the pulmonary artery to pass to the aorta.
Changes at birth
The ductus arteriosus stays open because of circulating factors including prostaglandins. The foramen ovale stays open because of the flow of blood from the right atrium to the left atrium. As the lungs expand, blood flows easily through the lungs and the membranous portion of the foramen ovale (the septum primum) flops over the muscular portion (the septum secundum). If the closure is incomplete, the result is a patent foramen ovale. The two flaps may fuse, but many adults have a foramen ovale that stays closed only because of the pressure difference between the atria.
Theories
Rokitansky (1875) explained congenital heart defects as breaks in heart development at various ontogenesis stages. Spitzer (1923) treats them as returns to one of the phylogenesis stages. Krimski (1963), synthesizing two previous points of view, considered congenital heart diseases as a stop of development at the certain stage of ontogenesis, corresponding to this or that stage of the phylogenesis. Hence, these theories can explain feminine and neutral types of defects only.
Diagnosis
Many congenital heart defects can be diagnosed prenatally by fetal echocardiography. This is a test which can be done during the second trimester of pregnancy, when the woman is about 18–24 weeks pregnant. It can be an abdominal ultrasound or transvaginal ultrasound.
If a baby is born with cyanotic heart disease, the diagnosis is usually made shortly after birth due to the blue colour of their skin (called cyanosis).
If a baby is born with a septal defect or an obstruction defect, often their symptoms are only noticeable after several months, or sometimes even after many years.
Classification
A number of classification systems exist for congenital heart defects. In 2000 the International Congenital Heart Surgery Nomenclature was developed to provide a generic classification system.
Hypoplasia
Hypoplasia can affect the heart, typically resulting in the underdevelopment of the right ventricle or the left ventricle. This causes only one side of the heart to be capable of pumping blood to the body and lungs effectively. Hypoplasia of the heart is rare but is the most serious form of CHD. It is called hypoplastic left heart syndrome when it affects the left side of the heart and hypoplastic right heart syndrome when it affects the right side of the heart. In both conditions, the presence of a patent ductus arteriosus (and, when hypoplasia affects the right side of the heart, a patent foramen ovale) is vital to the infant's ability to survive until emergency heart surgery can be performed, since without these pathways blood cannot circulate to the body (or lungs, depending on which side of the heart is defective). Hypoplasia of the heart is generally a cyanotic heart defect.
Obstructive defects
Obstructive defects occur when heart valves, arteries, or veins are abnormally narrow or blocked. Common defects include pulmonic stenosis, aortic stenosis, and coarctation of the aorta, with other types such as bicuspid aortic valve stenosis and subaortic stenosis being comparatively rare. Any narrowing or blockage can cause heart enlargement or hypertension.
Septal defects
The septum is a wall of tissue which separates the left heart from the right heart. Defects in the interatrial septum or the interventricular septum allow blood to flow from the left side of the heart to the right, reducing the heart's efficiency. Ventricular septal defects are collectively the most common type of CHD, although approximately 30% of adults have a type of atrial septal defect called probe patent foramen ovale.
Cyanotic defects
Cyanotic heart defects are called such because they result in cyanosis, a bluish-grey discoloration of the skin due to a lack of oxygen in the body. Such defects include persistent truncus arteriosus, total anomalous pulmonary venous connection, tetralogy of Fallot, transposition of the great vessels, and tricuspid atresia.
Defects
Aortic stenosis
Arrhythmogenic right ventricular cardiomyopathy
Atrial septal defect (ASD)
Atrioventricular septal defect (AVSD)
Bicuspid aortic valve
Cardiomyopathy
Complete heart block (CHB)
Dextrocardia
Double inlet left ventricle (DILV)
Double outlet right ventricle (DORV)
Ebstein's anomaly
Early Repolarization Syndrome
Holmes heart
Hypoplastic left heart syndrome (HLHS)
Hypoplastic right heart syndrome (HRHS)
Mitral stenosis
Myocardial bridge
Persistent truncus arteriosus
Pulmonary atresia
Pulmonary stenosis
Rhabdomyomas (Tumors of the Heart)
Transposition of the great vessels
dextro-Transposition of the great arteries (d-TGA)
levo-Transposition of the great arteries (l-TGA)
Tricuspid atresia
Ventricular septal defect (VSD)
Wolff–Parkinson–White syndrome (WPW)
Some conditions affect the great vessels or other vessels in close proximity to the heart, but not the heart itself, but are often classified as congenital heart defects.
Coarctation of the aorta (CoA)
Double aortic arch, aberrant subclavian artery, and other malformations of the great arteries
Interrupted aortic arch (IAA)
Patent ductus arteriosus (PDA)
Scimitar syndrome (SS)
Partial anomalous pulmonary venous connection (PAPVC)
Total anomalous pulmonary venous connection (TAPVC)
Some constellations of multiple defects are commonly found together.
Tetralogy of Fallot (ToF)
Pentalogy of Cantrell
Shone's syndrome/ Shone's complex / Shone's anomaly
Treatment
CHD may require surgery and medications. Medications include diuretics, which aid the body in eliminating water, salts, and digoxin for strengthening the contraction of the heart. This slows the heartbeat and removes some fluid from tissues. Some defects require surgical procedures to restore circulation back to normal and in some cases, multiple surgeries are needed.
Interventional cardiology now offers minimally invasive alternatives to surgery for some patients. The Melody Transcatheter Pulmonary Valve (TPV), approved in Europe in 2006 and in the U.S. in 2010 under a Humanitarian Device Exemption (HDE), is designed to treat congenital heart disease patients with a dysfunctional conduit in their right ventricular outflow tract (RVOT). The RVOT is the connection between the heart and lungs; once blood reaches the lungs, it is enriched with oxygen before being pumped to the rest of the body. Transcatheter pulmonary valve technology provides a less-invasive means to extend the life of a failed RVOT conduit and is designed to allow physicians to deliver a replacement pulmonary valve via a catheter through the patient's blood vessels.
Many people require lifelong specialized cardiac care, first with a pediatric cardiologist and later with an adult congenital cardiologist. There are more than 1.8 million adults living with congenital heart defects.
Mental health
Supporting people with chronic diseases such as congenital heart disease with emotional problems and mental health is a treatment consideration. Since some people with congenital heart disease have a lower quality of life that is related to their condition, some people may struggle with finding a job, engaging in physical exercise, with their fertility, and clinical depression as examples. An estimated 31% of adults with congenital heart disease also have mood disorders. Psychotherapy may be helpful for treating some people who have congenital heart disease and depression, however further research is needed to determine the best way to reduce depression including the length of treatments required for an improvement, type of psychotherapy treatments, and how the psychotherapy sessions are delivered.
Epidemiology
Heart defects are among the most common birth defect, occurring in 1% of live births (2–3% including bicuspid aortic valve). In 2013, 34.3 million people had CHD. In 2010, they resulted in 223,000 deaths, down from 278,000 deaths in 1990.
For congenital heart defects that arise without a family history (de novo), the recurrence risk in offspring is 3–5%. This risk is higher in left ventricular outflow tract obstructions, heterotaxy, and atrioventricular septal defects.
Terminology
Congenital heart defects are known by a number of names including congenital heart anomaly, congenital heart disease, heart defects, and congenital cardiovascular malformations.
| Biology and health sciences | Cardiovascular disease | Health |
868983 | https://en.wikipedia.org/wiki/Stretching | Stretching | Stretching is a form of physical exercise in which a specific muscle or tendon (or muscle group) is deliberately expanded and flexed in order to improve the muscle's felt elasticity and achieve comfortable muscle tone. The result is a feeling of increased muscle control, flexibility, and range of motion. Stretching is also used therapeutically to alleviate cramps and to improve function in daily activities by increasing range of motion.
In its most basic form, stretching is a natural and instinctive activity; it is performed by humans and many other animals. It can be accompanied by yawning. Stretching often occurs instinctively after waking from sleep, after long periods of inactivity, or after exiting confined spaces and areas. In addition to vertebrates (e.g. mammals and birds), spiders have also been found to exhibit stretching.
Increasing flexibility through stretching is one of the basic tenets of physical fitness. It is common for athletes to stretch before (for warming up) and after exercise in an attempt to reduce risk of injury and increase performance.
Stretching can be dangerous when performed incorrectly. There are many techniques for stretching in general, but depending on which muscle group is being stretched, some techniques may be ineffective or detrimental, even to the point of causing hypermobility, instability, or permanent damage to the tendons, ligaments, and muscle fiber. The physiological nature of stretching and theories about the effect of various techniques are therefore subject to heavy inquiry.
Although static stretching is part of some warm-up routines, pre-exercise static stretching usually reduces an individual's overall muscular strength and maximal performance, regardless of an individual's age, sex, or training status. For this reason, an active dynamic warm-up is recommended before exercise in place of static stretching.
Physiology
Studies have shed light on the function, in stretching, of a large protein within the myofibrils of skeletal muscles named titin. A study performed by Magid and Law demonstrated that the origin of passive muscle tension (which occurs during stretching) is actually within the myofibrils, not extracellularly as had previously been supposed. Due to neurological safeguards against injury such as the Golgi tendon reflex, it is normally impossible for adults to stretch most muscle groups to their fullest length without training due to the activation of muscle antagonists as the muscle reaches the limit of its normal range of motion.
Psychology
Stretching has been recognized for its potential to positively influence both cognitive function and mood. Research indicates that engaging in stretching exercises may lead to a reduction in feelings of anxiety, depression, hostility, fatigue, and confusion, particularly among individuals with sedentary lifestyles. These improvements in mood have been observed to correlate with enhancements in cognitive function. For individuals who often spend prolonged periods engaged in sedentary activities, integrating stretching into their daily routines may prove beneficial. Doing so not only addresses physical tension but also promotes mental well-being. Regular stretching has been associated with decreased levels of anxiety and depression, alongside increased vigor, which could activate brain regions associated with improved cognitive abilities.
Types of stretches
Stretches can be either static or dynamic. Static stretches are performed while stationary and dynamic stretches involve movement of the muscle. Stretches can also be active or passive, where active stretches use internal forces generated by the body to perform a stretch and passive stretches involve forces from external objects or people to perform the stretch. They can involve both passive and active components.
Dynamic stretching
Dynamic stretching is a movement-based stretch aimed at increasing blood flow throughout the body while also loosening up the muscle fibers. Standard dynamic stretches typically involve slow and controlled active contraction of muscles. An example of such a dynamic stretch is lunges. Another form of dynamic stretching is ballistic stretching, which is an active stretch that involves bouncing or swinging back and forth at a high speed in order to take a muscle beyond its typical range of motion using momentum. Ballistic stretching can also be performed with tools such as resistance bands to increase the intention between sets in order to quickly warm-up the body. Ballistic stretching may cause damage to the joints.
Static stretching
The simplest static stretches are static–passive stretches, according to research findings. This brings the joint to its end range of motion and hold it there using external forces. There are more advanced forms of static stretching, such as proprioceptive neuromuscular facilitation (PNF), which involves both active muscle contractions and passive external forces. PNF stretching utilizes an aspect of neuromuscular reeducation, which may yield better results than regular static stretching in terms of induced strength. PNF stretching may involve contracting either the antagonist muscles, agonist muscles, or both (CRAC). The efficacy of PNF stretching and its recommendation of use may be dependent on stretching-to-performance duration.
Effectiveness
Stretching has been found both effective and ineffective based on its application for treatment.
Although many people engage in stretching before or after exercise, the medical evidence has shown this has no meaningful benefit in preventing specifically muscle soreness. It may reduce the lactic acid build up in the muscles, making the next workout more bearable.
Stretching does not appear to reduce the risk of injury during exercises, except perhaps a dynamic warm-up for runners. While running places extreme stress loads on the joints, static stretching can help to improve joint flexibility. However, this has not been proven to reduce risk of injury in the runners. A dynamic (stretching) warm up has been shown to help overall running performance.
Delayed onset muscle soreness, also known as DOMS, typically arises 48 hours after an exercise bout. Stretching before or after the exercise did not show any significant benefits in the onset of DOMS.
Effectiveness of dynamic stretching
Ballistic stretching, a form of dynamic stretching, is likely to increase flexibility through a neurological mechanism. The stretched muscle is moved passively to the end range by an external force or agonist muscle: holding a muscle in this position might reduce muscle spindle sensitivity, with repeated stretch applied at the end range inhibiting the Golgi tendon organ.
Dynamic stretching, because it is movement-based, may not isolate the muscle group as well or have as intense of a stretch, but it is better at increasing the circulation of blood flow throughout the body, which in turn increases the amount of oxygen able to be used for athletic performance. This type of stretching has shown better results on athletic performances of power and speed when compared to static stretching.
Effectiveness of static stretching
Static stretching is better at creating a more intense stretch because it is able to isolate a muscle group better. But this intensity of stretching may hinder one's athletic performance because the muscle is being overstretched while held in this position and, once the tension is released, the muscle will tend to tighten up and may actually become weaker than it was previously. It has been shown in high level athletes, such as gymnasts, after performing a static stretching routine that it has a negative effect. The gymnasts lost the ability to jump vertically as high as prior as well as no improvement in their straddle jump or flexibility. Also, the longer the duration of static stretching, the more exhausted the muscle becomes. This type of stretching has been shown to have negative results on athletic performance within the categories of power and speed. However, to be able to do usual daily activities, a certain amount of range of motion is needed from each muscle. For example, the calf muscles are one of the muscle groups that have the most need for adequate flexibility since they are deeply related to normal lower limb function. When the goal is to increase flexibility, the most commonly used technique is stretching. Chronic static stretching was shown to increase range of motion of Dorsiflexion or bringing one's foot closer to their shin by an average of 5.17 degrees in healthy individuals versus 3.77 degrees when solely using ballistic stretching. While static stretching is shown to decrease power and speed in higher level athletes, when it comes to the older population who live more sedentary lifestyles static stretching has been shown to increase muscles strength and power.
Dynamic versus Static: Flexibility and Performance
Both dynamic and static stretching have been shown to have a positive impact on flexibility over time by increasing muscle and joint elasticity, thus increasing the depth and range of motion an athlete is able to reach. This is evident in the experiment "Acute effects of duration on sprint performance of adolescent football players". In this experiment, football players were put through different stretching durations of static and dynamic stretching to test their effects. They were tested on maximum sprinting ability and overall change in flexibility. Both static and dynamic stretching had a positive impact on flexibility but, whereas dynamic stretching had no impact on sprint times, static stretching had a negative result, worsening the time the participants were able to sprint the distance in.
| Biology and health sciences | Physical fitness | Health |
869123 | https://en.wikipedia.org/wiki/Gastroenteritis | Gastroenteritis | Gastroenteritis, also known as infectious diarrhea, is an inflammation of the gastrointestinal tract including the stomach and intestine. Symptoms may include diarrhea, vomiting, and abdominal pain. Fever, lack of energy, and dehydration may also occur. This typically lasts less than two weeks. Although it is not related to influenza, in the U.S. and U.K., it is sometimes called the "stomach flu".
Gastroenteritis is usually caused by viruses; however, gut bacteria, parasites, and fungi can also cause gastroenteritis. In children, rotavirus is the most common cause of severe disease. In adults, norovirus and Campylobacter are common causes. Eating improperly prepared food, drinking contaminated water or close contact with a person who is infected can spread the disease. Treatment is generally the same with or without a definitive diagnosis, so testing to confirm is usually not needed.
For young children in impoverished countries, prevention includes hand washing with soap, drinking clean water, breastfeeding babies instead of using formula, and proper disposal of human waste. The rotavirus vaccine is recommended as a prevention for children. Treatment involves getting enough fluids. For mild or moderate cases, this can typically be achieved by drinking oral rehydration solution (a combination of water, salts and sugar). In those who are breastfed, continued breastfeeding is recommended. For more severe cases, intravenous fluids may be needed. Fluids may also be given by a nasogastric tube. Zinc supplementation is recommended in children. Antibiotics are generally not needed. However, antibiotics are recommended for young children with a fever and bloody diarrhea.
In 2015, there were two billion cases of gastroenteritis, resulting in 1.3 million deaths globally. Children and those in the developing world are affected the most. In 2011, there were about 1.7 billion cases, resulting in about 700,000 deaths of children under the age of five. In the developing world, children less than two years of age frequently get six or more infections a year. It is less common in adults, partly due to the development of immunity.
Signs and symptoms
Gastroenteritis usually involves both diarrhea and vomiting. Sometimes, only one or the other is present. This may be accompanied by abdominal cramps. Signs and symptoms usually begin 12–72 hours after contracting the infectious agent. If due to a virus, the condition usually resolves within one week. Some viral infections also involve fever, fatigue, headache and muscle pain. If the stool is bloody, the cause is less likely to be viral and more likely to be bacterial. Some bacterial infections cause severe abdominal pain and may persist for several weeks.
Children infected with rotavirus usually make a full recovery within three to eight days. However, in poor countries treatment for severe infections is often out of reach and persistent diarrhea is common. Dehydration is a common complication of diarrhea. Severe dehydration in children may be recognized if the skin color and position returns slowly when pressed. This is called "prolonged capillary refill" and "poor skin turgor". Abnormal breathing is another sign of severe dehydration. Repeat infections are typically seen in areas with poor sanitation, and malnutrition. Stunted growth and long-term cognitive delays can result.
Reactive arthritis occurs in 1% of people following infections with Campylobacter species. Guillain–Barré syndrome occurs in 0.1%. Hemolytic uremic syndrome (HUS) may occur due to infection with Shiga toxin-producing Escherichia coli or Shigella species. HUS causes low platelet counts, poor kidney function, and low red blood cell count (due to their breakdown). Children are more predisposed to getting HUS than adults. Some viral infections may produce benign infantile seizures.
Cause
Viruses (particularly rotavirus (in children) and norovirus (in adults)) and the bacteria Escherichia coli and Campylobacter species are the primary causes of gastroenteritis. There are, however, many other infectious agents that can cause this syndrome including parasites and fungi. Non-infectious causes are seen on occasion, but they are less likely than a viral or bacterial cause. Risk of infection is higher in children due to their lack of immunity. Children are also at higher risk because they are less likely to practice good hygiene habits. Children living in areas without easy access to water and soap are especially vulnerable.
Viral
Rotaviruses, noroviruses, adenoviruses, and astroviruses are known to cause viral gastroenteritis. Rotavirus is the most common cause of gastroenteritis in children, and produces similar rates in both the developed and developing world. Viruses cause about 70% of episodes of infectious diarrhea in the pediatric age group. Rotavirus is a less common cause in adults due to acquired immunity. Norovirus is the cause in about 18% of all cases. Generally speaking, viral gastroenteritis accounts for 21–40% of the cases of infectious diarrhea in developed countries.
Norovirus is the leading cause of gastroenteritis among adults in America accounting for about 90% of viral gastroenteritis outbreaks. These localized epidemics typically occur when groups of people spend time proximate to each other, such as on cruise ships, in hospitals, or in restaurants. People may remain infectious even after their diarrhea has ended. Norovirus is the cause of about 10% of cases in children.
Bacterial
In some countries, Campylobacter jejuni is the primary cause of bacterial gastroenteritis, with half of these cases associated with exposure to poultry. In children, bacteria are the cause in about 15% of cases, with the most common types being Escherichia coli, Salmonella, Shigella, and Campylobacter species. If food becomes contaminated with bacteria and remains at room temperature for several hours, the bacteria multiply and increase the risk of infection in those who consume the food. Some foods commonly associated with illness include raw or undercooked meat, poultry, seafood, and eggs; raw sprouts; unpasteurized milk and soft cheeses; and fruit and vegetable juices. In the developing world, especially sub-Saharan Africa and Asia, cholera is a common cause of gastroenteritis. This infection is usually transmitted by contaminated water or food.
Toxigenic Clostridioides difficile is an important cause of diarrhea that occurs more often in the elderly. Infants can carry these bacteria without developing symptoms. It is a common cause of diarrhea in those who are hospitalized and is frequently associated with antibiotic use. Staphylococcus aureus infectious diarrhea may also occur in those who have used antibiotics. Acute "traveler's diarrhea" is usually a type of bacterial gastroenteritis, while the persistent form is usually parasitic. Acid-suppressing medication appears to increase the risk of significant infection after exposure to several organisms, including Clostridioides difficile, Salmonella, and Campylobacter species. The risk is greater in those taking proton pump inhibitors than with H2 antagonists.
Parasitic
A number of parasites can cause gastroenteritis. Giardia lamblia is most common, but Entamoeba histolytica, Cryptosporidium spp., and other species have also been implicated. As a group, these agents comprise about 10% of cases in children. Giardia occurs more commonly in the developing world, but this type of illness can occur nearly everywhere. It occurs more commonly in persons who have traveled to areas with high prevalence, children who attend day care, men who have sex with men, and following disasters.
Transmission
Transmission may occur from drinking contaminated water or when people share personal objects. Water quality typically worsens during the rainy season and outbreaks are more common at this time. In areas with four seasons, infections are more common in the winter. Worldwide, bottle-feeding of babies with improperly sanitized bottles is a significant cause. Transmission rates are also related to poor hygiene, (especially among children), in crowded households, and in those with poor nutritional status. Adults who have developed immunities might still carry certain organisms without exhibiting symptoms. Thus, adults can become natural reservoirs of certain diseases. While some agents (such as Shigella) only occur in primates, others (such as Giardia) may occur in a wide variety of animals.
Non-infectious
There are a number of non-infectious causes of inflammation of the gastrointestinal tract. Some of the more common include medications (like NSAIDs), certain foods such as lactose (in those who are intolerant), and gluten (in those with celiac disease). Crohn's disease is also a non-infectious cause of (often severe) gastroenteritis. Disease secondary to toxins may also occur. Some food-related conditions associated with nausea, vomiting, and diarrhea include: ciguatera poisoning due to consumption of contaminated predatory fish, scombroid associated with the consumption of certain types of spoiled fish, tetrodotoxin poisoning from the consumption of puffer fish among others, and botulism typically due to improperly preserved food.
In the United States, rates of emergency department use for noninfectious gastroenteritis dropped 30% from 2006 until 2011. Of the twenty most common conditions seen in the emergency department, rates of noninfectious gastroenteritis had the largest decrease in visits in that time period.
Pathophysiology
Gastroenteritis is defined as vomiting or diarrhea due to inflammation of the small or large bowel, often due to infection. The changes in the small bowel are typically noninflammatory, while the ones in the large bowel are inflammatory. The number of pathogens required to cause an infection varies from as few as one (for Cryptosporidium) to as many as 108 (for Vibrio cholerae).
Diagnosis
Gastroenteritis is typically diagnosed clinically, based on a person's signs and symptoms. Determining the exact cause is usually not needed as it does not alter the management of the condition.
However, stool cultures should be performed in those with blood in the stool, those who might have been exposed to food poisoning, and those who have recently traveled to the developing world. It may also be appropriate in children younger than 5, old people, and those with poor immune function. Diagnostic testing may also be done for surveillance. As hypoglycemia occurs in approximately 10% of infants and young children, measuring serum glucose in this population is recommended. Electrolytes and kidney function should also be checked when there is a concern about severe dehydration.
Dehydration
A determination of whether or not the person has dehydration is an important part of the assessment, with dehydration typically divided into mild (3–5%), moderate (6–9%), and severe (≥10%) cases. In children, the most accurate signs of moderate or severe dehydration are a prolonged capillary refill, poor skin turgor, and abnormal breathing. Other useful findings (when used in combination) include sunken eyes, decreased activity, a lack of tears, and a dry mouth. A normal urinary output and oral fluid intake is reassuring. Laboratory testing is of little clinical benefit in determining the degree of dehydration. Thus the use of urine testing or ultrasounds is generally not needed.
Differential diagnosis
Other potential causes of signs and symptoms that mimic those seen in gastroenteritis that need to be ruled out include appendicitis, volvulus, inflammatory bowel disease, urinary tract infections, and diabetes mellitus. Pancreatic insufficiency, short bowel syndrome, Whipple's disease, coeliac disease, and laxative abuse should also be considered. The differential diagnosis can be complicated somewhat if the person exhibits only vomiting or diarrhea (rather than both).
Appendicitis may present with vomiting, abdominal pain, and a small amount of diarrhea in up to 33% of cases. This is in contrast to the large amount of diarrhea that is typical of gastroenteritis. Infections of the lungs or urinary tract in children may also cause vomiting or diarrhea. Classical diabetic ketoacidosis (DKA) presents with abdominal pain, nausea, and vomiting, but without diarrhea. One study found that 17% of children with DKA were initially diagnosed as having gastroenteritis.
Prevention
Water, sanitation, hygiene
A supply of easily accessible uncontaminated water and good sanitation practices are important for reducing rates of infection and clinically significant gastroenteritis. Personal hygiene measures (such as hand washing with soap) have been found to decrease rates of gastroenteritis in both the developing and developed world by as much as 30%. Alcohol-based gels may also be effective. Food or drink that is thought to be contaminated should be avoided.
Breastfeeding is important, especially in places with poor hygiene, as is improvement of hygiene generally. Breast milk reduces both the frequency of infections and their duration.
Vaccination
Due to both its effectiveness and safety, in 2009 the World Health Organization recommended that the rotavirus vaccine be offered to all children globally. Two commercial rotavirus vaccines exist and several more are in development. In Africa and Asia these vaccines reduced severe disease among infants and countries that have put in place national immunization programs have seen a decline in the rates and severity of disease. This vaccine may also prevent illness in non-vaccinated children by reducing the number of circulating infections. Since 2000, the implementation of a rotavirus vaccination program in the United States has substantially decreased the number of cases of diarrhea by as much as 80 percent. The first dose of vaccine should be given to infants between 6 and 15 weeks of age. The oral cholera vaccine has been found to be 50–60% effective over two years.
There are a number of vaccines against gastroenteritis in development. For example, vaccines against Shigella and enterotoxigenic Escherichia coli (ETEC), which are two of the leading bacterial causes of gastroenteritis worldwide.
Management
Gastroenteritis is usually an acute and self-limiting disease that does not require medication. The preferred treatment in those with mild to moderate dehydration is oral rehydration therapy (ORT). For children at risk of dehydration from vomiting, taking a single dose of the anti vomiting medication metoclopramide or ondansetron, may be helpful, and butylscopolamine is useful in treating abdominal pain.
Rehydration
The primary treatment of gastroenteritis in both children and adults is rehydration. This is preferably achieved by drinking rehydration solution, although intravenous delivery may be required if there is a decreased level of consciousness or if dehydration is severe. Drinking replacement therapy products made with complex carbohydrates (i.e. those made from wheat or rice) may be superior to those based on simple sugars. Drinks especially high in simple sugars, such as soft drinks and fruit juices, are not recommended in children under five years of age as they may increase diarrhea. Plain water may be used if more specific ORT preparations are unavailable or the person is not willing to drink them. A nasogastric tube can be used in young children to administer fluids if warranted. In those who require intravenous fluids, one to four hours' worth is often sufficient.
Dietary
It is recommended that breast-fed infants continue to be nursed in the usual fashion, and that formula-fed infants continue their formula immediately after rehydration with ORT. Lactose-free or lactose-reduced formulas usually are not necessary. Children should continue their usual diet during episodes of diarrhea with the exception that foods high in simple sugars should be avoided. The BRAT diet (bananas, rice, applesauce, toast and tea) is no longer recommended, as it contains insufficient nutrients and has no benefit over normal feeding.
A Cochrane Review from 2020 concludes that probiotics make little or no difference to people who have diarrhea lasting 2 days or longer and that there is no proof that they reduce its duration. They may be useful in preventing and treating antibiotic associated diarrhea. Fermented milk products (such as yogurt) are similarly beneficial. Zinc supplementation appears to be effective in both treating and preventing diarrhea among children in the developing world.
Antiemetics
Antiemetic medications may be helpful for treating vomiting in children. Ondansetron has some utility, with a single dose being associated with less need for intravenous fluids, fewer hospitalizations, and decreased vomiting. Metoclopramide might also be helpful. However, the use of ondansetron might possibly be linked to an increased rate of return to hospital in children. The intravenous preparation of ondansetron may be given orally if clinical judgment warrants. Dimenhydrinate, while reducing vomiting, does not appear to have a significant clinical benefit.
Antibiotics
Antibiotics are not usually used for gastroenteritis, although they are sometimes recommended if symptoms are particularly severe or if a susceptible bacterial cause is isolated or suspected. If antibiotics are to be employed, a macrolide (such as azithromycin) is preferred over a fluoroquinolone due to higher rates of resistance to the latter. Pseudomembranous colitis, usually caused by antibiotic use, is managed by discontinuing the causative agent and treating it with either metronidazole or vancomycin. Bacteria and protozoans that are amenable to treatment include Shigella Salmonella typhi, and Giardia species. In those with Giardia species or Entamoeba histolytica, tinidazole treatment is recommended and superior to metronidazole. The World Health Organization (WHO) recommends the use of antibiotics in young children who have both bloody diarrhea and fever.
Antimotility agents
Antimotility medication has a theoretical risk of causing complications, and although clinical experience has shown this to be unlikely, these drugs are discouraged in people with bloody diarrhea or diarrhea that is complicated by fever. Loperamide, an opioid analogue, is commonly used for the symptomatic treatment of diarrhea. Loperamide is not recommended in children, however, as it may cross the immature blood–brain barrier and cause toxicity. Bismuth subsalicylate, an insoluble complex of trivalent bismuth and salicylate, can be used in mild to moderate cases, but salicylate toxicity is theoretically possible.
Epidemiology
It is estimated that there were two billion cases of gastroenteritis that resulted in 1.3 million deaths globally in 2015. Children and those in the developing world are most commonly affected. As of 2011, in those younger than five, there were about 1.7 billion cases resulting in 0.7 million deaths, with most of these occurring in the world's poorest nations. More than 450,000 of these fatalities are due to rotavirus in children under five years of age. Cholera causes about three to five million cases of disease and kills approximately 100,000 people yearly. In the developing world, children less than two years of age frequently get six or more infections a year that result in significant gastroenteritis. It is less common in adults, partly due to the development of acquired immunity.
In 1980, gastroenteritis from all causes caused 4.6 million deaths in children, with the majority occurring in the developing world. Death rates were reduced significantly (to approximately 1.5 million deaths annually) by 2000, largely due to the introduction and widespread use of oral rehydration therapy. In the US, infections causing gastroenteritis are the second most common infection (after the common cold), and they result in between 200 and 375 million cases of acute diarrhea and approximately ten thousand deaths annually, with 150 to 300 of these deaths in children less than five years of age.
Society and culture
Gastroenteritis is associated with many colloquial names, including "Montezuma's revenge", "Delhi belly", "la turista", and "back door sprint", among others. It has played a role in many military campaigns and is believed to be the origin of the term "no guts no glory".
Gastroenteritis is the main reason for 3.7 million visits to physicians a year in the United States and 3 million visits in France. In the United States gastroenteritis as a whole is believed to result in costs of US$23 billion per year, with rotavirus alone resulting in estimated costs of US$1 billion a year.
Terminology
The first usage of "gastroenteritis" was in 1825. Before this time it was commonly known as typhoid fever or "cholera morbus", among others, or less specifically as "griping of the guts", "surfeit", "flux", "colic", "bowel complaint", or any one of several other archaic names for acute diarrhea. Cholera morbus is a historical term that was used to refer to gastroenteritis rather than specifically cholera.
Animals
Many of the same agents cause gastroenteritis in cats and dogs as in humans. The most common organisms are Campylobacter, Clostridioides difficile, Clostridium perfringens, and Salmonella. A large number of toxic plants may also cause symptoms.
Some agents are more specific to a certain species. Transmissible gastroenteritis coronavirus (TGEV) occurs in pigs resulting in vomiting, diarrhea, and dehydration. It is believed to be introduced to pigs by wild birds and there is no specific treatment available. It is not transmissible to humans.
| Biology and health sciences | Illness and injury | null |
869797 | https://en.wikipedia.org/wiki/Neurospora%20crassa | Neurospora crassa | Neurospora crassa is a type of red bread mold of the phylum Ascomycota. The genus name, meaning 'nerve spore' in Greek, refers to the characteristic striations on the spores. The first published account of this fungus was from an infestation of French bakeries in 1843.
Neurospora crassa is used as a model organism because it is easy to grow and has a haploid life cycle that makes genetic analysis simple since recessive traits will show up in the offspring. Analysis of genetic recombination is facilitated by the ordered arrangement of the products of meiosis in Neurospora ascospores. Its entire genome of seven chromosomes has been sequenced.
Neurospora was used by Edward Tatum and George Wells Beadle in their experiments for which they won the Nobel Prize in Physiology or Medicine in 1958. Beadle and Tatum exposed N. crassa to x-rays, causing mutations. They then observed failures in metabolic pathways caused by errors in specific enzymes. This led them to propose the "one gene, one enzyme" hypothesis that specific genes code for specific proteins. Their hypothesis was later elaborated to enzyme pathways by Norman Horowitz, also working on Neurospora. As Norman Horowitz reminisced in 2004, "These experiments founded the science of what Beadle and Tatum called 'biochemical genetics'. In actuality, they proved to be the opening gun in what became molecular genetics and all developments that have followed from that."
In the 24 April 2003 issue of Nature, the genome of N. crassa was reported as completely sequenced. The genome is about 43 megabases long and includes approximately 10,000 genes. There is a project underway to produce strains containing knockout mutants of every N. crassa gene.
In its natural environment, N. crassa lives mainly in tropical and sub-tropical regions. It can be found growing on dead plant matter after fires.
Neurospora is actively used in research around the world. It is important in the elucidation of molecular events involved in circadian rhythms, epigenetics and gene silencing, cell polarity, cell fusion, development, as well as many aspects of cell biology and biochemistry.
The sexual cycle
Sexual fruiting bodies (perithecia) can only be formed when two mycelia of different mating type come together (see Figure). Like other Ascomycetes, N. crassa has two mating types that, in this case, are symbolized by A and a. There is no evident morphological difference between the A and a mating type strains. Both can form abundant protoperithecia, the female reproductive structure (see Figure). Protoperithecia are formed most readily in the laboratory when growth occurs on solid (agar) synthetic medium with a relatively low source of nitrogen. Nitrogen starvation appears to be necessary for expression of genes involved in sexual development. The protoperithecium consists of an ascogonium, a coiled multicellular hypha that is enclosed in a knot-like aggregation of hyphae. A branched system of slender hyphae, called the trichogyne, extends from the tip of the ascogonium projecting beyond the sheathing hyphae into the air. The sexual cycle is initiated (i.e. fertilization occurs) when a cell (usually a conidium) of opposite mating type contacts a part of the trichogyne (see Figure). Such contact can be followed by cell fusion leading to one or more nuclei from the fertilizing cell migrating down the trichogyne into the ascogonium. Since both A and a strains have the same sexual structures, neither strain can be regarded as exclusively male or female. However, as a recipient, the protoperithecium of both the A and a strains can be thought of as the female structure, and the fertilizing conidium can be thought of as the male participant.
The subsequent steps following fusion of A and a haploid cells, have been outlined by Fincham and Day and Wagner and Mitchell. After fusion of the cells, the further fusion of their nuclei is delayed. Instead, a nucleus from the fertilizing cell and a nucleus from the ascogonium become associated and begin to divide synchronously. The products of these nuclear divisions (still in pairs of unlike mating type, i.e. A/a) migrate into numerous ascogenous hyphae, which then begin to grow out of the ascogonium. Each of these ascogenous hypha bends to form a hook (or crozier) at its tip and the A and a pair of haploid nuclei within the crozier divide synchronously. Next, septa form to divide the crozier into three cells. The central cell in the curve of the hook contains one A and one a nucleus (see Figure). This binuclear cell initiates ascus formation and is called an "ascus-initial" cell. Next the two uninucleate cells on either side of the first ascus-forming cell fuse with each other to form a binucleate cell that can grow to form a further crozier that can then form its own ascus-initial cell. This process can then be repeated multiple times.
After formation of the ascus-initial cell, the A and a nucleus fuse with each other to form a diploid nucleus (see Figure). This nucleus is the only diploid nucleus in the entire life cycle of N. crassa. The diploid nucleus has 14 chromosomes formed from the two fused haploid nuclei that had 7 chromosomes each. Formation of the diploid nucleus is immediately followed by meiosis. The two sequential divisions of meiosis lead to four haploid nuclei, two of the A mating type and two of the a mating type. One further mitotic division leads to four A and four a nucleus in each ascus. Meiosis is an essential part of the life cycle of all sexually reproducing organisms, and in its main features, meiosis in N. crassa seems typical of meiosis generally.
As the above events are occurring, the mycelial sheath that had enveloped the ascogonium develops as the wall of the perithecium, becomes impregnated with melanin, and blackens. The mature perithecium has a flask-shaped structure.
A mature perithecium may contain as many as 300 asci, each derived from identical fusion diploid nuclei. Ordinarily, in nature, when the perithecia mature the ascospores are ejected rather violently into the air. These ascospores are heat resistant and, in the lab, require heating at 60 °C for 30 minutes to induce germination. For normal strains, the entire sexual cycle takes 10 to 15 days. In a mature ascus containing eight ascospores, pairs of adjacent spores are identical in genetic constitution, since the last division is mitotic, and since the ascospores are contained in the ascus sac that holds them in a definite order determined by the direction of nuclear segregations during meiosis. Since the four primary products are also arranged in sequence, a first division segregation pattern of genetic markers can be distinguished from a second division segregation pattern.
Fine structure genetic analysis
Because of the above features N. crassa was found to be very useful for the study of genetic events occurring in individual meioses. Mature asci from a perithecium can be separated on a microscope slide and the spores experimentally manipulated. These studies usually involved the separate culture of individual ascospores resulting from a single meiotic event and determining the genotype of each spore. Studies of this type, carried out in several different laboratories, established the phenomenon of "gene conversion" (e.g. see references).
As an example of the gene conversion phenomenon, consider genetic crosses of two N. crassa mutant strains defective in gene pan-2. This gene is necessary for the synthesis of pantothenic acid (vitamin B5), and mutants defective in this gene can be experimentally identified by their requirement for pantothenic acid in their growth medium. The two pan-2 mutations B5 and B3 are located at different sites in the pan-2 gene, so that a cross of B5 ´ B3 yields wild-type recombinants at low frequency. An analysis of 939 asci in which the genotypes of all meiotic products (ascospores) could be determined found 11 asci with an exceptional segregation pattern. These included six asci in which there was one wild-type meiotic product but no expected reciprocal double-mutant (B5B3) product. Furthermore, in three asci the ratio of meiotic products was 1B5:3B3, rather than in the expected 2:2 ratio. This study, as well as numerous additional studies in N. crassa and other fungi (reviewed by Whitehouse), led to an extensive characterization of gene conversion. It became clear from this work that gene conversion events arise when a molecular recombination event happens to occur near the genetic markers under study (e.g. pan-2 mutations in the above example). Thus studies of gene conversion allowed insight into the details of the molecular mechanism of recombination. Over the decades since the original observations of Mary Mitchell in 1955, a sequence of molecular models of recombination have been proposed based on both emerging genetic data from gene conversion studies and studies of the reaction capabilities of DNA. Current understanding of the molecular mechanism of recombination is discussed in the Wikipedia articles Gene conversion and Genetic recombination. An understanding of recombination is relevant to several fundamental biologic problems, such the role of recombination and recombinational repair in cancer (see BRCA1) and the adaptive function of meiosis (see Meiosis).
Adaptive function of mating type
That mating in N. crassa can only occur between strains of different mating types suggests that some degree of outcrossing is favored by natural selection. In haploid multicellular fungi, such as N. crassa, meiosis occurring in the brief diploid stage is one of their most complex processes. Although physically much larger than the diploid stage, the haploid multicellular vegetative stage characteristically has a simple modular construction with little differentiation. In N. crassa, recessive mutations affecting the diploid stage of the life cycle are quite frequent in natural populations. These mutations, when homozygous in the diploid stage, often cause spores to have maturation defects or to produce barren fruiting bodies with few ascospores (sexual spores). Most of these homozygous mutations cause abnormal meiosis (e.g., disturbed chromosome pairing or pachytene or diplotene). The number of genes affecting the diploid stage was estimated to be at least 435 (about 4% of the total number of 9,730 genes). Thus, outcrossing, promoted by the necessity for the union of opposite mating types, likely provides the benefit of masking recessive mutations that would otherwise be harmful to sexual spore formation (see Complementation (genetics)).
Current research
Neurospora crassa is not only a model organism for the study of phenotypic types in knock-out variants, but a particularly useful organism widely used in computational biology and the circadian clock. It has a natural reproductive cycle of 22 hours and is influenced by external factors such as light and temperature. Knock out variants of wild type N. crassa are widely studied to determine the influence of particular genes (see Frequency (gene)).
| Biology and health sciences | Basics | Plants |
869868 | https://en.wikipedia.org/wiki/White%20spirit | White spirit | White spirit (AU, UK and Ireland) or mineral spirits (US, Canada), also known as mineral turpentine (AU/NZ/ZA), turpentine substitute, and petroleum spirits, is a petroleum-derived clear liquid used as a common organic solvent in painting. There are also terms for specific kinds of white spirit, including Stoddard solvent and solvent naphtha (petroleum). White spirit is often used as a paint thinner, or as a component thereof, though paint thinner is a broader category of solvent. Odorless mineral spirits (OMS) have been refined to remove the more toxic aromatic compounds, and are recommended for applications such as oil painting.
A mixture of aliphatic, open-chain or alicyclic C7 to C12 hydrocarbons, white spirit is insoluble in water and is used as an extraction solvent, as a cleaning solvent, as a degreasing solvent and as a solvent in aerosols, paints, wood preservatives, lacquers, varnishes, and asphalt products. In western Europe about 60% of the total white spirit consumption is used in paints, lacquers and varnishes. White spirit is the most widely used solvent in the paint industry. In households, white spirit is commonly used to clean paint brushes after use, to clean auto parts and tools, as a starting fluid for charcoal grills, to remove adhesive residue from non-porous surfaces, and many other common tasks.
The word "mineral" in "mineral spirits" or "mineral turpentine" is meant to distinguish it from distilled spirits (alcoholic beverages distilled from fermented biological material) or from true turpentine (distilled tree resin, composed mostly of pinene).
Types and grades
Three different types and three different grades of white spirit exist. The type refers to whether the solvent has been subjected to hydrodesulfurization (removal of sulfur) alone (type 1), solvent extraction (type 2) or hydrogenation (type 3).
Each type comprises three grades: low flash grade, regular grade, and high flash grade (flash refers to flash point). The grade is determined by the crude oil used as the starting material and the conditions of distillation.
In addition there is type 0, which is defined as distillation fraction with no further treatment, consisting predominantly of saturated C9 to C12 hydrocarbons with a boiling range of .
Stoddard solvent is a specific mixture of hydrocarbons, typically over 65% C10 or higher hydrocarbons, developed in 1924 by Atlanta dry cleaner W. J. Stoddard and Lloyd E. Jackson of the Mellon Institute of Industrial Research as a less flammable petroleum-based dry cleaning solvent than the petroleum solvents then in use. Dry cleaners began using the result of their work in 1928 and it soon became the predominant dry cleaning solvent in the United States, until the late 1950s.
Turpentine substitute is generally not made to a standard and can have a wider range of components than products marketed as white spirit, which is made to a standard (in the UK, British Standard BS 245, in Germany, DIN 51632). Turpentine substitute can be used for general cleaning but is not recommended for paint thinning as it may adversely affect drying times due to the less volatile components; while it may be used for brush cleaning its heavier components may leave an oily residue.
Chemical registry numbers
Physical properties
Type 1 white spirit is mainly used in most of Europe and Stoddard solvent is used in the US, both of which correspond to each other.
Use
Degreasing and lubricating
In industry, white spirit is used for cleaning and degreasing machine tools and parts, and in conjunction with cutting oil as a thread cutting and reaming lubricant.
White spirit is commonly used for cutting fluid in ultraprecision lathes (commonly referred to as diamond turning machines).
White spirit is used for regripping golf clubs. After the old grip is removed, the white spirit is poured into the new grip and shaken. After the white spirit is poured on, the new underlying tape and the new grip are slid on. After an hour of drying out, the new grip and club are ready to use.
Solvent and paint thinner
White Spirit is a petroleum distillate used as a paint thinner and mild solvent.
White spirit is an inexpensive petroleum-based replacement for the vegetable-based turpentine. It is commonly used as a paint thinner for oil-based paint and cleaning brushes, and as an organic solvent in other applications. Mineral turpentine is chemically very different from turpentine, which mainly consists of pinene, and it has inferior solvent properties. Artists use white spirit as an alternative to turpentine since it is less flammable and less toxic. Because of interactions with pigments in oil paints, artists require a higher grade of white spirit than many industrial users, including the complete absence of residual sulfur.
White spirit was formerly an active ingredient in the laundry soap Fels Naptha, used to dissolve oils and grease in laundry stains, and as a popular remedy for eliminating the irritant oil urushiol in poison ivy. It was removed as a potential health risk.
White spirit has a characteristic unpleasant kerosene-like odor. Chemical manufacturers have developed a low odor version of mineral turpentine which contains less of the highly volatile shorter hydrocarbons. Odorless mineral spirits is white spirit that has been further refined to remove the more toxic aromatic compounds, and is recommended for applications such as oil painting, where humans have close contact with the solvent.
In screen printing (also referred to as silk-screening), white spirit is often used to clean and unclog screens after printing with oil-based textile and plastisol inks. It is also used to thin inks used in making monoprints.
White spirit is often used inside liquid-filled compasses and gauges.
White spirits are a major ingredient in some popular automotive fuel/oil additives, such as Marvel Mystery Oil, as they are capable of dissolving varnish and sludge buildup.
Portable lanterns and stoves
Although white spirit is sometimes used as an alternative to camp fuel, such as kerosene or paraffin, in portable lanterns and camp stoves , this is highly inadvisable as typical grades of white spirit have a lower flash point than kerosene. It cannot be used as an alternative to Coleman camp fuel or white gas, which is a much more volatile gasoline-like fuel.
Toxicity
White spirit is mainly classed as an irritant. It has a fairly low acute toxicity by inhalation of the vapour, dermal (touching the skin) and oral (ingestion) routes. However, acute exposure can lead to central nervous system depression resulting in lack of coordination and slowed reactions. Exposure to very high concentrations in enclosed spaces can lead to general narcotic effects (drowsiness, dizziness, nausea, etc.) and can eventually lead to unconsciousness. Oral ingestion presents a high aspiration hazard. Prolonged or repeated skin exposure over a long period of time can result in severe irritant dermatitis, also called contact dermatitis.
Continuous exposure to an average white spirit concentration of 240 mg/m3 (40 ppm) for more than 13 years can lead to chronic central nervous system effects. Similar long-term studies have been made in which some of the observed effects included memory impairment, poor concentration, increased irritability etc. White spirit is implicated in the development of chronic toxic encephalopathy (CTE) among house painters. In severe cases CTE may lead to disability and personality changes. These effects in painters were first studied in the 1970s in the Nordic countries.
Owing to the volatility and low bioavailability of its constituents, white spirit, although it is moderately toxic to aquatic organisms, is unlikely to present significant hazards to the environment. It should not however, be purposely poured down the sink or freshwater drain.
People can be exposed to Stoddard solvent in the workplace by breathing it in, swallowing it, skin contact, and eye contact. The Occupational Safety and Health Administration (OSHA) has set the legal limit (Permissible exposure limit) for Stoddard solvent exposure in the workplace as 500 ppm (2900 mg/m3) over an 8-hour workday. The National Institute for Occupational Safety and Health (NIOSH) has set a recommended exposure limit (REL) of 350 mg/m3 over an 8-hour workday and 1800 mg/m3 over 15 minutes. At levels of 20,000 mg/m3, Stoddard solvent is immediately dangerous to life and health.
| Physical sciences | Hydrocarbons | Chemistry |
869944 | https://en.wikipedia.org/wiki/Geological%20formation | Geological formation | A geological formation, or simply formation, is a body of rock having a consistent set of physical characteristics (lithology) that distinguishes it from adjacent bodies of rock, and which occupies a particular position in the layers of rock exposed in a geographical region (the stratigraphic column). It is the fundamental unit of lithostratigraphy, the study of strata or rock layers.
A formation must be large enough that it can be mapped at the surface or traced in the subsurface. Formations are otherwise not defined by the thickness of their rock strata, which can vary widely. They are usually, but not universally, tabular in form. They may consist of a single lithology (rock type), or of alternating beds of two or more lithologies, or even a heterogeneous mixture of lithologies, so long as this distinguishes them from adjacent bodies of rock.
The concept of a geologic formation goes back to the beginnings of modern scientific geology. The term was used by Abraham Gottlob Werner in his theory of the origin of the Earth, which was developed over the period from 1774 to his death in 1817. The concept became increasingly formalized over time and is now codified in such works as the North American Stratigraphic Code and its counterparts in other regions.
Geologic maps showing where various formations are exposed at the surface are fundamental to such fields as structural geology, allowing geologists to infer the tectonic history of a region or predict likely locations for buried mineral resources.
Defining formations
The boundaries of a formation are chosen to give it the greatest practical lithological consistency. Formations should not be defined by any criteria other than lithology. The lithology of a formation includes characteristics such as chemical and mineralogical composition, texture, color, primary depositional structures, fossils regarded as rock-forming particles, or other organic materials such as coal or kerogen. The taxonomy of fossils is not a valid lithological basis for defining a formation.
The contrast in lithology between formations required to justify their establishment varies with the complexity of the geology of a region. Formations must be able to be delineated at the scale of geologic mapping normally practiced in the region; the thickness of formations may range from less than a meter to several thousand meters.
Geologic formations are typically named after a permanent natural or artificial feature of the geographic area in which they were first described. The name consists of the geographic name plus either "Formation" or a descriptive name. Examples include the Morrison Formation, named for the town of Morrison, Colorado, and the Kaibab Limestone, named after the Kaibab Plateau of Arizona. The names must not duplicate previous formation names, so, for example, a newly designated formation could not be named the Kaibab Formation, since the Kaibab Limestone is already established as a formation name. The first use of a name has precedence over all others, as does the first name applied to a particular formation.
As with other stratigraphic units, the formal designation of a formation includes a stratotype which is usually a type section. A type section is ideally a good exposure of the formation that shows its entire thickness. If the formation is nowhere entirely exposed, or if it shows considerably lateral variation, additional reference sections may be defined. Long-established formations dating to before the modern codification of stratigraphy, or which lack tabular form (such as volcanic formations), may substitute a type locality for a type section as their stratotype. The geologist defining the formation is expected to describe the stratotype in sufficient detail that other geologists can unequivocally recognize the formation.
Although formations should not be defined by any criteria other than primary lithology, it is often useful to define biostratigraphic units on paleontological criteria, chronostratigraphic units on the age of the rocks, and chemostratigraphic units on geochemical criteria, and these are included in stratigraphic codes.
Usefulness of formations
The concept of formally defined layers or strata is central to the geologic discipline of stratigraphy, and the formation is the fundamental unit of stratigraphy. Formations may be combined into groups of strata or divided into members. Members differ from formations in that they need not be mappable at the same scale as formations, though they must be lithologically distinctive where present.
The definition and recognition of formations allow geologists to correlate geologic strata across wide distances between outcrops and exposures of rock strata. Formations were at first described as the essential geologic time markers, based on their relative ages and the law of superposition. The divisions of the geological time scale were described and put in chronological order by the geologists and stratigraphers of the 18th and 19th centuries.
Geologic formations can be usefully defined for sedimentary rock layers, low-grade metamorphic rocks, and volcanic rocks. Intrusive igneous rocks and highly metamorphosed rocks are generally not considered to be formations, but are described instead as lithodemes.
Other uses of the term
"Formation" is also used informally to describe the odd shapes (forms) that rocks acquire through erosional or depositional processes. Such a formation is abandoned when it is no longer affected by the geologic agent that produced it. Some well-known cave formations include stalactites and stalagmites.
| Physical sciences | Stratigraphy | Earth science |
1334840 | https://en.wikipedia.org/wiki/Beluga%20%28sturgeon%29 | Beluga (sturgeon) | The beluga (), also known as the beluga sturgeon or great sturgeon (Huso huso), is a species of anadromous fish in the sturgeon family (Acipenseridae) of the order Acipenseriformes. It is found primarily in the Caspian and Black Sea basins, and formerly in the Adriatic Sea. Based on maximum size, it is the third-most-massive living species of bony fish. Heavily fished for the female's valuable roe, known as beluga caviar, wild populations have been greatly reduced by overfishing and poaching, leading IUCN to classify the species as critically endangered.
Etymology
The common name for the sturgeon, as for the unrelated beluga whale, is derived from the Russian word (), meaning , probably referring to the extensive pale colour on the flanks and belly in beluga compared to that of other sturgeons.
Description
Huso huso shows typical characteristics of other sturgeon, such as an elongated body, heterocercal tail, partially cartilaginous skeleton, naked skin and longitudinal series of scutes.
The dorsal fin has 48 to 81 soft rays, and the anal fin, much shorter, has 22 to 41 soft rays. There are five in a series of longitudinal scutes: dorsal (one series, 9–17 scutes), lateral (two series, one per side, 28–60 scutes each) and ventral series (two series, one per side, 7–14 scutes each). The surface of the skin is covered by fine denticles. The rostrum is conical and contains numerous sensory pits on both ventral and dorsal surfaces. The mouth is large, crescent-shaped and protractile, with the upper lip continuous while the lower lip is interrupted by a large gap. The barbels are laterally compressed with foliate appendages, arranged in two pairs, originating midway or closer to the mouth than to the tip of the snout.
However, during growth, the beluga sturgeons show evident morphologic changes:
Juveniles are slender, and the head is quite narrow with a mouth ventrally placed but projecting upward. The snout is thin and pointed (almost half of the head), scutes are evident, back and flancs are dark grey or black and the belly is white.
Adults are heavy-set, spindle-shaped, large and humpbacked. The head is massive with a very large protractile mouth that gradually moves in an almost frontal position during growth. The snout is quite short (one-third to one-quarter of the head), and scutes gradually undergo absorption and decrease in number with age. Colouring is blue-grey or dark brown, with silver or grey flanks and white belly. The dark dorsum contrasts strongly against the rest of the body;
Very old specimens are stocky, with a large head and an enormous mouth.
Size
Among all extant bony fishes, the beluga sturgeon rivals the ocean sunfish (Mola sp.) as the most massive fish and is the second-longest bony fish after the giant oarfish (Regalecus glesne). It is the largest freshwater fish in the world. The beluga also rivals the great white shark (Carcharodon carcharias) and the greenland shark (Somniosus microcephalus) for the title of largest actively predatory fish.
The largest accepted record is of a female taken in 1827 in the Volga estuary at and . Another specimen reportedly weighed and measured in length. Claims about greater length (, or even ); and weight (, , or even ) are disputed and unconfirmed; but they are not impossible. Several other records of aged sturgeon exceed . Among sturgeons, only the closely related Kaluga (Huso dauricus) can attain similar size, with a maximum weight of .
Beluga of such great sizes are very old (continuing to grow throughout life) and have become increasingly rare in recent decades because of heavy fishing of the species. Today, mature belugas that are caught are generally long and weigh . The female beluga is typically 20% larger than the male. An exceptionally large beluga recently caught weighed and measured .
Biology
Due to rampant overfishing, the average lifespan of beluga sturgeon is unknown, with no specimens living past their 56th year. However, the species is reportedly quite long-lived, being capable of surviving over 100 years in the wild.
Spawning
Like most sturgeons, the beluga is anadromous, migrating upstream in rivers to spawn on clean, hard substrate, which offers both support and cover to their sticky and adhesive eggs. Spawning biology and development of larval stages of the sturgeon, the most ancient fish of the Danube, co-evolved with the formation of the Danube valley, resulting in very different survival strategies in its early life stages. This appears to explain why different individuals of the same long-migratory species spawn as far upstream as upstream, while others spawn just . To make the long journey to very distant spawning grounds, the sturgeon adapted a two-stage migration strategy, beginning in autumn when they enter the Danube River overwinter in the river and the second stage is their spawning which takes place in the spring the river in fall and staying over winter in reaches of the river offering adequate substrate and water-flow resting conditions. Very few locations of existing wintering and spawning grounds for sturgeon are presently known in the lower Danube, and none are known to exist in the river's upper reaches. The same situation concerns nursery sites upon which young sturgeon depend during their journey to the Black Sea.
Males attain sexual maturity at 12–16 years of age, whereas females do so at 16–22 years. They will spawn every four to seven years. At one time, beluga sturgeons could migrate up to upriver to spawn, but dams in almost every major tributary that they utilize have impeded historic spawning routes. The female lays her eggs on gravel from underwater. Upon hatching, the embryo are long, and 10–14 days later when they absorb their yolk sack, the length is . While swimming back to the ocean, the young sturgeon may cover up to a day.
Diet
Huso huso is a pelagic predator whose local distribution is not influenced by the nature of the substrates, unlike with most of the sturgeons that show demersal attitude. The prey is sucked into the mouth opening extremely quickly.
Juveniles feed on benthic invertebrates in rivers and shallow coastal waters, where they grow quickly. At the length of , they become largely piscivorous.
Different diets have been observed throughout the distribution range of beluga sturgeon, as well as according to spawners' migration stage. Adults mainly eat a great diversity of large fish (73% of the diet). Additional food items may include molluscs and crustaceans, aquatic birds and young seals (Caspian seals, Pusa caspica).
The piscivorous diet of beluga sturgeon tends to change with age: in the Caspian Sea, it mainly consists of Clupeonella sp. for juveniles smaller than 40 cm, different species of Gobiidae for fish ranging between 40 and 280 cm and then mullets, Alosa sp. and other sturgeons for the largest.
In brackish environments of the Ponto-Caspian basin, the genera Alosa, Aspius and Engraulis are the preferred prey. In estuaries and rivers of the same area, migrating spawners eat various cyprinids, mainly Cyprinus carpio and Rutilus rutilus, Sander lucioperca and, among sturgeons, Acipenser ruthenus is the main prey.
Little is known about the diet of the extinct Adriatic population. It has been reported that in marine and brackish environments, adult Adriatic H. huso foraged primarily on molluscs (Cephalopoda, of which common cuttlefish, Sepia officinalis, and European squid, Loligo vulgaris, are particularly common in the Adriatic Sea) and fish belonging to the families Gadidae, Pleuronectidae, Gobiidae, Clupeidae, Scombridae and Mugilidae, but also on big crustaceans; in the rivers, they fed mainly on local Cyprinidae.
Habitat
Beluga sturgeon are considered euryhaline, capable of moving freely between freshwater and estuaries, and thus can live in waters of varying saline content. Sturgeons are quite a durable species and can survive some of the most altered and polluted rivers in the world.
Historically, beluga sturgeon were found in the Caspian Sea, Black Sea, Adriatic Sea, Sea of Azov, and all rivers interconnecting these waterways. Unfortunately, this range has been greatly reduced in modern times to the Caspian Sea, the Black Sea, and a few rivers such as the Danube, with attempts to reintroduce Belugas into various historic locations.
Uses
Beluga caviar is considered a delicacy worldwide. The flesh of the beluga is not particularly renowned, but it is a hearty white meat similar to that of swordfish. Beluga caviar has long been scarce and expensive and the fish's endangered status has made its caviar even more expensive throughout the world.
The beluga's air bladder is said to make the best isinglass.
Status
IUCN classifies the beluga as critically endangered. Due to aforementioned poaching of the sturgeon, the Danube is the only river remaining with naturally reproducing sturgeon populations within the European Union. The sturgeon remains a protected species listed in Appendix III of the Bern Convention, and its trade is restricted under CITES Appendix II. The Mediterranean population is strongly protected under Appendix II of the Bern Convention, prohibiting any intentional killing of these fish.
The United States Fish and Wildlife Service has banned imports of beluga caviar and other beluga products from the Caspian Sea since 6 October 2005, after listing beluga sturgeon under the U.S. Endangered Species Act.
Repopulation efforts
Since 2015, an official captive breeding scheme has been established in Italy, with beluga from the Azov Sea. Then, after the building of a fish ladder on Isola Serafini dam, on 2019 hundreds of young microchipped beluga and 60 tagged subadults were released into the Po river, following EU Projects (Life Ticino Biosurce). Since then, many H. huso were released in the Po river, attempting to resurrect the extinct Adriatic population.
Management of sturgeon fisheries within the Caspian Sea began in the 1950s and while the initial regulations had honorable intentions, they achieved dwindling effects due to the ever present demand for the fish's caviar. In July 2016, Sturgeon Aquafarms in Bascom, Florida, became the first and only facility in the world to obtain a permit exemption for the sale of beluga sturgeon and its caviar in the U.S. Since 2017, the company has assisted in beluga sturgeon repopulation efforts across the world by providing over 160,000 fertilized eggs to the Caspian Sea region.
Following a World Wildlife Fund crowdfunding appeal, over 7,000 three month-old beluga sturgeons were released into the Danube River. Despite repopulation efforts, the beluga sturgeon continues to face poaching threats. In 2021, two Romanian men in Grindu, Ialomita, were caught trying to smuggle a 140 kilogram, 2.5-metre beluga sturgeon in a wagon; the fish was later safely returned to the river.
Threats
The beluga sturgeon is confronted by several critical threats that imperil its existence. Illegal harvesting, habitat disruption through dam construction, and pollution are among the most pressing challenges faced.
Illegal harvesting and poaching
The beluga sturgeon faces a significant threat from illegal catches for its meat and caviar. The excessive harvesting and a sharp increase in poaching have removed the largest and most mature specimens from the population, almost eliminating natural reproduction. This exploitation has pushed the species to the brink of extinction.
Habitat impoundment
The construction of dams, such as the Iron Gate in the Danube and the Volgograd Dam, has severely reduced the beluga sturgeon's available spawning grounds. These dams have blocked access to crucial river habitats, leading to a significant reduction in the species' ability to reproduce. Similar habitat loss has occurred in other rivers due to dam construction, greatly impacting the species' survival.
Pollution
Pollution from various sources, including oil, industries, sewage, and agriculture, is a critical threat to the beluga sturgeon. The species' long lifespan makes it vulnerable to pesticide contamination, resulting in reduced reproductive success and other health issues. Pollution negatively affects the quality of the sturgeon's habitat, compounding the challenges faced by this endangered species.
| Biology and health sciences | Acipenseriformes | Animals |
1335268 | https://en.wikipedia.org/wiki/American%20badger | American badger | The American badger (Taxidea taxus) is a North American badger similar in appearance to the European badger, although not closely related. It is found in the western, central, and northeastern United States, northern Mexico, and south-central Canada to certain areas of southwestern British Columbia.
The American badger's habitat is typified by open grasslands with available prey (such as mice, squirrels, and groundhogs). The species prefers areas such as prairie regions with sandy loam soils where it can dig more easily for its prey.
Taxonomy
The American badger is a member of the Mustelidae, a diverse family of carnivorous mammals that also includes weasels, otters, ferrets, and the wolverine. The American badger belongs to the Taxidiinae, one of four subfamilies of mustelid badgers – the other three being the Melinae (four species in two genera, including the European badger), the Helictidinae (five species of ferret-badgers) and the Mellivorinae (the honey badger); the so-called stink badgers are mephitids. The American badger's closest relative is the prehistoric Chamitataxus. Among extant mustelids, the American badger is the most basal species; its lineage is thought to have split off from the rest of the Mustelidae about 18 million years (Ma) ago, following the split of mustelids from procyonids about 29 Ma ago.
The recognized subspecies include:
The ranges of the subspecies overlap considerably, with intermediate forms occurring in the areas of overlap.
In Mexico, this animal is sometimes called tlalcoyote. The Spanish word for badger is tejón, but in Mexico this word is also used to describe the coati. This can lead to confusion, as both coatis and badgers are found in Mexico.
Description
The American badger has most of the general characteristics common to badgers; with stocky and low-slung bodies with short, powerful legs, they are identifiable by their huge foreclaws (measuring up to 5 cm in length) and distinctive head markings.
American badgers possess morphological characteristics that enable them to be good fossorial specialists, such as a conical head, bristles on the ears, and nictitating membranes in the eyes. American badgers have powerful forelimbs. They also possess a strong humerus and large bony processes for the attachment of muscles. The mechanical advantage in badger forelimbs is increased by the specialized olecranon process and bones such as the radius and metacarpals.
Measuring generally between in length, males of the species are slightly larger than females. They may attain an average weight of roughly for females and up to for males. Northern subspecies such as T. t. jeffersonii are heavier than the southern subspecies. In the fall, when food is plentiful, adult male badgers can reach up to . In some northern populations, females can average .
Except for the head, the American badger is covered with a grizzled, brown, black and white coat of coarse hair or fur, giving almost a mixed brown-tan appearance. The coat aids in camouflage in grassland habitat. Its triangular face shows a distinctive black and white pattern, with brown or blackish "badges" marking the cheeks and a white stripe extending from the nose to the base of the head. In the subspecies T. t. berlandieri, the white head stripe extends the full length of the body, to the base of the tail.
Diet
The American badger is a fossorial carnivore. It preys predominantly on pocket gophers (Geomyidae), ground squirrels (Spermophilus), moles (Talpidae), marmots (Marmota), prairie dogs (Cynomys), pika (Ochotona), woodrats (Neotoma), kangaroo rats (Dipodomys), deer mice (Peromyscus), and voles (Microtus), often digging to pursue prey into their dens, and sometimes plugging tunnel entrances with objects. The American badger is a significant predator of snakes, including rattlesnakes, and is considered their most important predator in South Dakota. They also prey on ground-nesting birds, such as the bank swallow, or sand martin (Riparia riparia), and the burrowing owl (Athene cunicularia), lizards, amphibians, carrion, fish, skunks (Mephitis and Spilogale), insects (including bees and honeycomb), and some plant foods, such as corn (Zea mais), peas, green beans, mushrooms and other fungi, and sunflower seeds (Helianthus).
Behavior
American badgers are generally nocturnal; however, in remote areas with no human encroachment they are routinely observed foraging during the day. Seasonally, a badger observed during daylight hours in the spring months of late March to early May often represents a female foraging during daylight and spending nights with her young. Badgers do not hibernate but may become less active in winter. A badger may spend much of the winter in cycles of torpor that last around 29 hours. They do emerge from their burrows when the temperature is above freezing.
As a fossorial mammal, the American badger uses a scratch-digging process where the forelimbs are withdrawn to break the soil and move the debris behind or to the sides of its body
An abandoned badger burrow may be occupied by mammals of similar size, such as foxes and skunks, as well as animals as diverse as the burrowing owl, tiger salamander and California red-legged frog.
The American badger has been seen working with a coyote in tandem while hunting. Typically this pairing is one badger to one coyote; however, one study found about 9% of sightings included two coyotes to one badger, while 1% had one badger to three coyotes. Researchers have found that the coyote benefits by an increased catch rate of about 33%, and while it is difficult to see precisely how the badger benefits, the badger has been noted to spend more time underground and active. Badgers are also thought to expend less energy while hunting in burrows.
According to research, this partnership works due to the different hunting styles of the predators and how their prey reacts to them. A ground squirrel, upon spotting a coyote, will crawl into its hole to escape; while upon seeing a badger, the ground squirrel will climb out of its hole and use its speed to outrun the badger. Hunting in tandem raises the prey vulnerability and both predators win.
Life cycle
Badgers are normally solitary animals but are thought to expand their territories in the breeding season to seek out mates. Mating occurs in late summer and early fall, with some males breeding with more than one female. American badgers experience delayed implantation, with pregnancies suspended until December or as late as February. Young are born from late March to early April in litters ranging from one to five young, averaging about three.
Badgers are born blind, furred, and helpless. Eyes open at four to six weeks. The female feeds her young solid foods prior to complete weaning and for a few weeks thereafter. Young American badgers first emerge from the den on their own at five to six weeks old. Families usually break up and juveniles disperse from the end of June to August; young American badgers leave their mothers as early as late May or June. Juvenile dispersal movements are erratic.
Most female American badgers become pregnant for the first time after they are a year old. A minority of females four to five months old ovulate, and a few become pregnant. Males usually do not breed until their second year.
Large predators occasionally kill American badgers. The average longevity in the wild is 9–10 years, with a record of 14; a captive example lived at least 15 years and five months.
Habitat
American badgers prefer grasslands and open areas with grasslands, which can include parklands, farms, and treeless areas with friable soil and a supply of rodent prey. They may also be found in forest glades and meadows, marshes, brushy areas, hot deserts, and mountain meadows. They are sometimes found at elevations up to but are usually found in the Sonoran and Transition life zones (which are at elevations lower and warmer than those characterized by coniferous forests). In Arizona, they occur in desert scrub and semi-arid grasslands. In California, American badgers are primarily able to survive through a combination of open grasslands of agricultural lands, protected land trust and open space lands, and regional and state and national park lands with grassland habitat. Badgers are occasionally found in open chaparral (with less than 50% plant cover) and riparian zones. They are not usually found in mature chaparral. In Manitoba aspen parklands, American badger abundance was positively associated with the abundance of Richardson's ground squirrels (Spermophilus richardsonii).
In Ontario it primarily resides on the extreme southwestern portion of the province, restricted to the north shore of Lake Erie in open areas generally associated with agriculture and along woodland edges. There have been a few reports from the Bruce-Grey region.
American badger use of home range varies with season and sex. Different areas of the home range are used more frequently at different seasons and usually are related to prey availability. Males generally have larger home ranges than females. In a 1972 study, radiotransmitter-tagged American badgers had an average annual home range of . The home range of one female was in summer, in fall, and in winter. Lindzey reported average home ranges of . Estimated density of American badgers in Utah scrub-steppe was one per square mile (2.6 km2), with 10 dens in active or recent use.
, overdevelopment of American badger habitat had resulted in reduced range, decreased prey, and forced badgers into contact with humans when foraging between fragments. Direct observations in Sonoma County, documenting habitat and badger sightings and foraging, reflect various ranges within the fragmented habitat areas from less than 1/2 mile to approximately 4 miles. Within these areas, the availability of prey and a fresh water source are key factors for the preferred habitat areas and ability to survive. Identifying and conserving habitat areas where there is year-round activity, along with identified burrowing patterns and observations of female badger territory for birthing and raising young have become critical factors in survival of the species.
Plant communities
American badgers are most commonly found in treeless areas, including tallgrass and shortgrass prairies, grass-dominated meadows and fields within forested habitats, and shrub-steppe communities. In the Southwest, plant indicators of the Sonoran and Transition life zones (relatively low, dry elevations) commonly associated with American badgers include creosotebush
(Larrea tridentata), junipers (Juniperus spp.), gambel oak (Quercus gambelii), willows (Salix spp.), cottonwoods (Populus spp.), ponderosa pine (Pinus ponderosa), grasses, and sagebrushes (Artemisia spp.).
In Colorado in 1977, American badgers were common in grass–forb and ponderosa pine habitats. In Kansas, they are common in tallgrass prairie dominated by big bluestem (Andropogon gerardi), little bluestem (Schizachyrium scoparium), and Indian grass (Sorghastrum nutans). In Montana in 1990, badgers were present in Glacier National Park in fescue (Festuca spp.) grasslands. In Manitoba, they occur in grassland extensions within aspen (Populus spp.) parklands.
Cover requirements
American badgers require cover for sleep, concealment, protection from weather, and natal denning. They typically enlarge foraged out gopher or other prey holes, or other animal burrows. Their dens range from about 4 feet to 10 feet in depth and 4 feet to 6 feet in width. A female American Badger may create 2 to 4 burrows in proximity with a connecting tunnel for concealment and safety for her young. Displaced soil from digging out the burrow characteristically appears in front of the burrow entrance, and a view from a distance reveals a mound-like roof of the burrow, with the living and concealment space created underneath the raised-roof appearing mound.
During summer and autumn, badgers range more frequently, with mating season generally in November, and burrowing patterns reflect 1 to 3 burrows may be dug from foraged out prey holes in a day, used for a day to a week, and then abandoned, with possible returns later, and other small wildlife utilizing abandoned burrows in the interim. Where prey is particularly plentiful, they will reuse dens, especially in the fall, sometimes for a few days at a time. In winter, a single den may be used for most of the season. Natal dens are dug by the female and are used for extended periods, but litters may be moved, probably to allow the mother to forage in new areas close to the nursery. Natal dens are usually larger and more complex than diurnal dens.
Predation
While the American badger is an aggressive animal with few natural enemies, it is still vulnerable to other species in its habitat. Predation on American badger by golden eagles (Aquila chrysaetos), coyotes (Canis latrans) and bobcats (Lynx rufus) has been reported. Bears (Ursus spp.) and gray wolves (Canis lupus) occasionally kill American badgers, while cougars (Puma concolor), according to a 2019 study, apparently are the main predators of adults, hunting them much more frequently than the other carnivorans, with a documented case where the badger is one of the main prey of a radio-collared cougar.
American badgers are trapped by humans for their pelts. Their fur is used for shaving and painting brushes.
Conservation status
In May 2000, the Canadian Species at Risk Act listed each of the subspecies Taxidea taxus jacksoni and T. t. jeffersonii as an endangered species in Canada. The nominate subspecies was additionally deemed to be of "Special Concern" in November 2012. The subspecies T. t. jeffersonii, which resides in British Columbia, was subsequently divided into two populations, a western one in the Okanagan Valley-Cariboo region and an eastern one in East Kootenay, with each receiving an endangered listing. The Canadian population of T. t. jacksoni is isolated from other badgers in a small area of southwestern Ontario near the border with the United States south of the Niagara Peninsula, although a population may also exist in the northwest of the province.
The California Department of Fish and Game designated the American badger as a California species of special concern.
| Biology and health sciences | Mustelidae | Animals |
1335820 | https://en.wikipedia.org/wiki/Inhaler | Inhaler | An inhaler (puffer, asthma pump or allergy spray) is a medical device used for delivering medicines into the lungs through the work of a person's breathing. This allows medicines to be delivered to and absorbed in the lungs, which provides the ability for targeted medical treatment to this specific region of the body, as well as a reduction in the side effects of oral medications. There are a wide variety of inhalers, and they are commonly used to treat numerous medical conditions with asthma and chronic obstructive pulmonary disease (COPD) being among the most notable.
Some of the common types of inhalers include metered-dose inhalers, dry powder inhalers, soft mist inhalers, and nebulizers. Each device has advantages and disadvantages and can be selected based on individually specific patient needs, as well as age, pathological conditions, coordination, and lung function. Proper education on inhaler use is important to ensure that inhaled medication creates its proper effects in the lungs. Using a spacer can ensure that more medicine reaches the lungs, thus providing the most optimal treatment.
Medical uses
Inhalers are designed to deliver medication directly to the lungs through a person's own breathing. This may benefit a patient by providing medicines directly to areas of disease, allowing medication to take a greater effect on its intended target, and limit side effects of medications when administered locally. Inhalers are used in a variety of different medical conditions with diseases of the lungs and respiratory system being among the most common. Individuals with these diseases/conditions need medications designed to decrease airway inflammation and obstruction to allow for easier and comfortable breathing. Antibiotic medications have even been developed for inhalers to allow for direct delivery to areas of infection within the lungs. Two of the most common conditions that warrant inhaler therapy are asthma and chronic obstructive pulmonary disease.
Asthma
Asthma is a condition of intermittent airway obstruction due to inflammatory processes in the lungs. Inhaled medications are used to calm down the inflammation present in the lungs and allow for relief of the airway obstruction. Common inhaled medications used for treatment of asthma include long term inhalational steroidal anti-inflammatory drugs (most commonly inhaled corticosteroids, also called ICS) and fast-relieving bronchodilators such as salbutamol (known commonly as "Ventolin") and salmeterol. These medications allow for patients to have relief of airway obstruction symptoms and reduced inflammation. If some people are unable to use inhalers, non-steroidal anti-inflammatory drugs (NSAIDs) may be used, but with caution since they may cause immunological hypersensitivity to NSAIDs, resulting in respiratory-related symptoms such as bronchospasms, acute asthma exacerbation, and severe asthma morbidity.
Chronic obstructive pulmonary disease (COPD)
COPD is an obstructive lung disease due to long-term damage to the airways of the lungs. The long-term damage leads to the inability of the airways to open properly, causing airway obstruction. Inhaled medications allow patients to see improvement in symptoms and better function of daily living. Some commonly used inhaled medications in patient's with COPD are ipratroprium, salmeterol, and corticosteroids. Inhalers that combine two or three different medications including inhaled corticosteroids, long-active muscarinic medications (LAMA) and long acting beta2 agonists (LABA) for treating COPD may be associated with improvements in some quality of life variables and small improvements in lung function and respiratory symptoms, however, may also be associated with an increase in the risk of pneumonia.
Types of inhalers
Meter-dosed inhaler (MDI)
The most common type of inhaler is the pressurized metered-dose inhaler (MDI) which is made up of 3 standard components- a metal canister, plastic actuator, and a metering valve. The medication is typically stored in solution in a pressurized canister that contains a propellant or suspension. The MDI canister is attached to a plastic, hand-operated actuator. On activation, the metered-dose inhaler releases a fixed dose of medication in aerosol form through the actuator and into a patient's lungs. These devices require significant coordination as a person must discharge the medication at or near the same time that they inhale in order for the medication to be effective.
Dry powder inhaler (DPI)
Dry powder inhalers release a metered or device-measured dose of powdered medication that is inhaled through a DPI device. This device usually contains a chamber in which the powdered medication is deposited prior to each dosage. The powder can then be inhaled with a quick breath. This allows for medication to be delivered to the lungs without the need for use of propellant/suspension.
Soft mist inhaler (SMI)
Soft mist inhalers release a light mist containing medication without the need for a propellant/suspension. Upon pressing a button, the inhaler creates a mist of medication, allowing for inhalation into the lungs. SMIs suspend inhaled medications for roughly 1.2 seconds, which is longer than the average MDI inhaler suspension time period. This requires less coordination when using and may be helpful for young patients or patients that find the MDI inhalers difficult to use.
Nebulizer
Nebulizers are designed to deliver medications over an extended period of time over multiple breaths through a mouthpiece or face mask. They generate a continuous mist with aerosolized medication, allowing a patient to breathe normally and receive medications. They are commonly used in infants and toddlers requiring inhaled medications or in patients in the hospital who require inhaled medications.
Smart inhaler
The smart-inhaler is an inhaler that will automatically update an app with information that includes the time of day, air quality, and how many times it has been used through sensor technology on the device. The first smart-inhaler was approved in 2019 by the FDA, its purpose is to track patient use of the device and some other circumstantial factors that could affect the effectiveness of the dosage. This information is sent via Bluetooth to a mobile device app, and is later shared with their physician to determine what kind of things can trigger issues with asthma and other problems. This technology presents a great way to cut down on medical costs associated with asthma and also help patients better manage their condition with fewer emergencies.
The Teva ProAir Digihaler was the first FDA approved smart inhaler. It shows how effective the device is at aiding patients in using the proper dose amount for their asthma. In a study published by the European Respiratory Journal, the ProAir Digihaler accurately identified when patients were using their inhalers and whether they were effectively administering the dose in a 370 patient trial with the device. This study further gives an overview on the technology regarding applications and devices that help aid in the tracking and medication management for asthma and other lung conditions. Another study showed that smart inhalers accurately recorded all doses administered by patients with their technology, which signifies their importance in providing accurate dosage information to patients and their physicians.
Propellants
In 2009, the FDA banned the use of inhalers that use chlorofluorocarbons (CFC) as propellants. In their place, inhalers now use hydrofluoroalkane (HFA). HFA is not environmentally inert as it is a greenhouse gas but it does not affect the ozone layer. While some people with asthma and advocacy groups contend that HFA inhalers are not as effective, published clinical studies indicate CFC and HFA inhalers are equally effective in controlling asthma.
While the impact of CFCs from inhalers on the ozone layer had been minuscule (dwarfed by industrial processes using CFCs), the FDA in its interpretation of the Montreal Protocol mandated the switch in propellants. Patients expressed concern about the high price of the HFA inhalers as there were initially no generic versions, whereas generic CFC inhalers had been available.
Proper use
It is important to use proper techniques when administering medications through inhalers.
Proper use of inhalers often involves initial deep breathing (which involves mostly the diaphragm's movements), and then rapid breathing (which involves most of the muscles of respiration, such as external and internal intercostal muscles) during intake of one or more puffs from the inhalers.
Improper use of inhalers is very common, can lead to distribution of the medicine into the mouth or throat where it cannot create its desired effect and may cause harm. Education on the correct use of inhalers for delivery of medications is a commonly cited topic in medical studies and a great deal of thought has been put into how best to help people learn to use their inhalers effectively. Below is a description of proper inhaler technique for each different type of inhaler as well as a helpful video explaining what the text states.
Meter-dosed inhalers
The mouthpiece is removed and the inhaler is shaken for 5–10 seconds.
The inhaler is gripped with mouthpiece on the bottom and canister on top. A finger is placed on the canister to allow for delivery of medicine.
Deep inhalation is done until no more air can be taken into the lungs.
Deep exhalation is done until most of the air is out of the lungs.
Once deep exhalation is done, mouth is placed over mouthpiece.
As the next deep inhalation begins, the canister is pressed down to release the medicine into the lungs.
Slow deep breathing is continued and breath is held for 5–10 seconds, keeping the medicine in the lungs for a longer time period and preventing escape of aerosolized form of the medicine.
Complete exhalation is done again. If multiple puffs of the medicine have to be taken, steps 1–5 are repeated after waiting for 15–30 seconds.
Mouthpiece is replaced.
With spacer
Spacer is placed at the mouthpiece of a meter-dosed inhaler while keeping mouth at the end of the spacer. After pressing the canister of the inhaler, the medicine will linger inside the spacer, allowing for the user of the inhaler to directly absorb medicine into their lungs. Deep breathing is done to be ready for the delivery of the medicine to the lungs; this minimizes need for coordination of breathing with inhaler activation. Cleanage of the spacer regularly with warm soapy water is recommended.
Dry powder inhalers
Inhaler medication chamber is prepared (this will be different based on the type of inhaler but will involve preparing and opening the chamber with the medication)
The inhaler is held with the chamber pointing towards the patient and complete exhalation is done with their head turned away from the inhaler.
Mouth is placed over the chamber and a quick, deep breath is taken allowing medication to dispense in the lungs.
Breath is held for 5–10 seconds and then slow exhalation is done.
After waiting for a few minutes, steps 1-4 are repeated if another dose is needed.
Soft mist inhalers
The inhaler is primed by loading the cartridge and discharging the inhaler until a fine mist is visible (more explanation in the video).
Once complete exhalation is done, mouth is placed around the mouthpiece while leaving space for the small holes on the side of the mouthpiece.
Slow inhalation is done while simultaneously pressing the button to release the medication.
Breath is held for 5–10 seconds.
Slow exhalation is done and steps 1-4 are repeated if another dose of medication is required after waiting for a few minutes.
If inhaler is used everyday, the inhaler usually has to be primed the first time using a new cartridge, but it may need to be primed again if it has not been used in multiple days.
After use
If using inhaled corticosteroids, one should wash the mouth out directly after use of an inhaler. This helps to prevent mouth infections that can occur due to immunosuppressant effects of corticosteroids.
Nebulizer
Mouth is placed over mouthpiece or face mask is placed over nose and mouth
The nebulizer machine is turned on.
Normal breathing is done for 10-20 min (or time allotted for treatment).
Machine is turned off and face mask/mouthpiece remove is removed.
Price and availability
In the United States, pharmaceutical manufacturers use legal and regulatory strategies to keep inhaler prices artificially high. There has been little innovation in inhaler technology for decades the most recent drug to be approved by the FDA for treating asthma or COPD via a novel target of action was Ipratropium bromide in 1986. Since then, manufacturers have used small changes to drug delivery mechanisms, or have switched active ingredients from one inhaler device to another (a strategy known as a "device hop") to keep patents active. This has the effect of limiting competition, keeping inhalers expensive. Because of high prices, patients sometimes skip doses or give up using their inhalers.
History
The idea of directly delivering medication into the lungs was based on ancient traditional cures that involved the use of aromatic and medicinal vapors. These did not involve any special devices beyond the apparatus used for burning or heating to produce fumes. Early inhalation devices included one devised by John Mudge in 1778. It had a pewter mug with a hole allowing attachment of a flexible tube. Mudge used it for the treatment of coughs using opium. These devices evolved with modifications by Wolfe, Mackenzie (1872) and better mouth attachments such as by Beigel in 1866. Many of these early inhalers needed heat to vaporize the active chemical ingredient. The benefits of forced expiration and inspiration to treat asthma were noted by J. S. Monell in 1865. Chemicals used in inhalers included ammonia, chlorine, iodine, tar, balsams, turpentine camphor and numerous others in combinations. Julius Mount Bleyer used a variation in 1890 in New York.
In 1968, Robert Wexler of Abbott Laboratories developed the Analgizer, a disposable inhaler that allowed the self-administration of methoxyflurane vapor in air for analgesia. The Analgizer consisted of a polyethylene cylinder 5 inches long and 1 inch in diameter with a 1 inch long mouthpiece. The device contained a rolled wick of polypropylene felt which held 15 milliliters of methoxyflurane.
Because of the simplicity of the Analgizer and the pharmacological characteristics of methoxyflurane, it was easy for patients to self-administer the drug and rapidly achieve a level of conscious analgesia which could be maintained and adjusted as necessary over a period of time lasting from a few minutes to several hours. The 15 milliliter supply of methoxyflurane would typically last for two to three hours, during which time the user would often be partly amnesic to the sense of pain; the device could be refilled if necessary.
The Analgizer was found to be safe, effective, and simple to administer in obstetric patients during childbirth, as well as for patients with bone fractures and joint dislocations, and for dressing changes on burn patients. When used for labor analgesia, the Analgizer allows labor to progress normally and with no apparent adverse effect on Apgar scores. All vital signs remain normal in obstetric patients, newborns, and injured patients. The Analgizer was widely utilized for analgesia and sedation until the early 1970s, in a manner that foreshadowed the patient-controlled analgesia infusion pumps of today. The Analgizer inhaler was withdrawn in 1974, but use of methoxyflurane as a sedative and analgesic continues in Australia and New Zealand in the form of the Penthrox inhaler.
| Biology and health sciences | General concepts_2 | Health |
1335911 | https://en.wikipedia.org/wiki/Usnea | Usnea | Usnea is a genus of fruticose lichens in the large family Parmeliaceae. The genus, which currently contains roughly 130 species, was established by Michel Adanson in 1763. Species in the genus grow like leafless mini-shrubs or tassels anchored on bark or twigs. Members of the genus are commonly called old man's beard, beard lichen, or beard moss.
Members of the genus are similar to those of the genus Alectoria. A distinguishing test is that the branches of Usnea are somewhat elastic, but the branches of Alectoria snap cleanly off.
Systematics
The genus Usnea was circumscribed by Michel Adanson in 1763. He used the name designated by Johann Jacob Dillenius, whose earlier published description did not meet the rules of valid publication as established by the International Code of Nomenclature for algae, fungi, and plants. However, he did not specify a type specimen; the species Usnea florida, moved to the genus by Friedrich Heinrich Wiggers in 1780, has been designated as the lectotype. Since the establishment of the genus, hundreds of Usnea species have been described. A three-volume series by Józef Motyka published in 1936 and 1947 listed 451 species. By 2006, the genus contained more than 600 species, which made it one of the largest genera within the family Parmeliaceae. However, many former species are now regarded as morphological varieties and adaptations to local circumstances. The number of recognized species in Finland has decreased for this reason, for example, dropping from 34 in 1951 to 25 in 1963 and only 12 in 2000. In addition, some former Usnea species have been moved to other genera; for instance, Usnea longissima was renamed Dolichousnea longissima in 2004. By 2022, the overall number of species assigned to the genus had dropped to 355.
The name Usnea is probably derived from the Arabic word Ushnah, meaning moss or lichen, though it may also mean "rope-like". Based on a fossil Usnea found in Baltic amber, the genus is known to date back to at least the late Eocene, about 34 million years ago.
Description
Usnea lichens are fruticose. Structurally they are shrubby, often with many branches, and can be erect or . Some trailing species can grow to considerable size; strands of Usnea longissima, for example, may exceed in length. Colours vary depending on the species, from straw-coloured, yellow-green or pale green through green or greyish-green to reddish or variegated red and green. Unlike other similar-looking fruticose lichens, species in this genus have an elastic chord or axis running through the middle of the thallus that can be revealed by gently pulling a filament apart from either end.
Usnea looks very similar to the plant Spanish moss, so much so that the latter's Latin name is derived from it (Tillandsia usneoides, the 'Usnea-like Tillandsia').
Distribution and habitat
Usnea lichens are widely distributed in both the northern and southern hemisphere, in both temperate and tropical regions.
They appear in areas with low levels of air pollution. They can often be found on the ground in areas with where trees or branches have recently been cut, such as orchards (after pruning) and active logging areas.
Ecology
Usnea lichens reproduce via vegetative means through fragmentation, asexual means through soredia, or sexual means through ascogonium and spermatogonium. The growth rate of lichens in nature is slow, but can be increased in laboratory conditions.
Like other lichens, Usnea often grows on sick or dying trees due to the pre-existing loss of canopy leaves, allowing for greater photosynthesis by the lichen's algae; this leads some gardeners to mistakenly blame the lichen for the tree's leaf loss and illness.
Usnea is very sensitive to air pollution, especially sulfur dioxide. Under poor growing conditions, such as areas high in pollution, they may grow no larger than a few millimetres, if they survive at all. Where the air is unpolluted, they can grow to 10–20 cm long. It can sometimes be used as a bioindicator, because it tends to only grow in those regions where the air is clean, and of high quality.
Uses
By humans
Traditional medicines
According to Paul Bergner, Author of Medical Herbalism, "the usnic acid in Usnea is effective against gram-positive bacteria such as Streptococcus and Staphylococcus, making Usnea a valuable addition to herbal formulas for sore throats and skin infections. It is also effective against a bacterium that commonly causes pneumonia."
Bolivian traditional healers called the Kallawaya use Kaka sunka in decoction to cure lung problems. The lichen is macerated in alcohol and rubbed onto the body of those suffering from "nervous fragility".
Some believe that Usnea, in high concentrations, could possess some toxicity. The National Toxicology Program evaluated the issue, undertaking research involved feeding male and female rats and mice ground Usnea lichens containing usnic acid for three months at various concentrations. Rats suffered severe toxicity, with significant liver damage observed at various concentrations, while mice experienced liver toxicity, ovarian atrophy, and changes in reproductive cycles at higher doses. Additionally, both species showed weight loss at elevated exposure levels, and mice exhibited potential genetic damage after two weeks at high concentrations. A safe exposure level was established at 60 parts per million, below which no adverse effects were observed.
Dyes
Usnea species have been used to create yellow, orange, green, blue, and purple dyes for textiles. This wide variety of possible colors can be achieved due to variations in chemical composition depending on the species, locality, and race of a particular specimen. Specifically, Usnea can contain thamnolic, squamatic, barbatic, salazinic, and alectorialic acids, all of which can affect dye color.
Many indigenous peoples of Central and South America, including the Tarahumura and Mapuche people, have a history of dyeing with Usnea, generally to obtain orange and brown hues. The Tarahumura use them to dye wool blankets in brown and russet hues, and the Mapuche have used Usnea florida to obtain orange. There is also anthropological evidence that Usnea cocca sonca was historically used for dyeing in Peru, and a lichen called cuaxapaxtle was used near Mexico City.
Cosmetics
Usnea barbata has been used in cosmetic production for its antimicrobial and antifungal properties as a preservative and deodorant.
Firestarters
When dry, Usnea lichens are flammable and can be used as a fire starter.
Food
Some Usnea species have been used as food sources during times of scarcity. For example, people in Bosnia and Herzegovina ate Usnea barbata during the Bosnian War, particularly in the winter, when other plant material was not readily available. They ground it into powdery "flour" to make bread or ate it as mush.
By other organisms
The northern parula, a species of New World warbler which breeds in North America, uses Usnea lichens in the construction of its nest in some parts of its range. Where these lichens have declined due to air pollution, the bird has also vanished as a breeding species.
Species
Usnea acromelana
Usnea alboverrucata
Usnea amblyoclada
Usnea angulata
Usnea antarctica
Usnea aranea
Usnea articulata
Usnea aurantiaciparvula
Usnea austrocampestris – Falkland Islands
Usnea bismolliuscula
Usnea boomiana
Usnea brattiae
Usnea cavernosa
Usnea cedrosiana
Usnea ceratina
Usnea chaetophora
Usnea cirrosa
Usnea clerciana
Usnea confusa
Usnea cornuta
Usnea crenulata
Usnea crocata
Usnea cylindrica
Usnea diplotypus
Usnea effusa
Usnea elata
Usnea elixii
Usnea esperantiana
Usnea exigua
Usnea filipendula
Usnea firmula
Usnea flammea
Usnea flavocardia
Usnea flavorubescens
Usnea fleigiae
Usnea florida
Usnea floriformis
Usnea foveata
Usnea fragilescens
Usnea fulvoreagens
Usnea galapagona
Usnea geissleriana
Usnea glabrata
Usnea glabrescens
Usnea glauca
Usnea grandisora
Usnea grandispora
Usnea himantodes
Usnea hirta
Usnea inermis
Usnea intermedia
Usnea kalbiana
Usnea krogiana
Usnea lambii
Usnea lapponica
Usnea leana
Usnea lutii
Usnea macaronesica
Usnea maculata
Usnea marivelensis
Usnea mayrhoferi
Usnea mekista
Usnea messutiae
Usnea molliuscula
Usnea myrmaiacaina
Usnea neuropogonoides
Usnea nidifica
Usnea nidulifera
Usnea oncodeoides
Usnea oncodes
Usnea oreophila
Usnea pacificana
Usnea pallidocarpa
Usnea parafloridana
Usnea patriciana
Usnea pendulina
Usnea perplexans
Usnea poliothrix
Usnea praetervisa
Usnea propagulifera
Usnea pulvinata
Usnea pycnoclada
Usnea pygmoidea
Usnea quasirigida
Usnea ramulosissima
Usnea roseola
Usnea rubicunda
Usnea rubricornuta
Usnea rubriglabrata
Usnea rubrotincta
Usnea sanguinea
Usnea saxidilatata
Usnea scabrata
Usnea scabrida
Usnea silesiaca
Usnea sphacelata
Usnea subalpina
Usnea subaranea
Usnea subcapillaris
Usnea subcomplecta
Usnea subcornuta
Usnea subdasaea
Usnea subeciliata
Usnea subflammea
Usnea subflaveola
Usnea subfloridana
Usnea subglabrata
Usnea subparvula
Usnea subrubicunda
Usnea subscabrosa
Usnea tamborensis
Usnea taylorii
Usnea torulosa
Usnea trachycarpa
Usnea ushuaiensis
Usnea viktoriana
Usnea vrieseana
Usnea wasmuthii
Usnea xanthopoga
| Biology and health sciences | Lichens | Plants |
1337421 | https://en.wikipedia.org/wiki/NGC%207027 | NGC 7027 | NGC 7027, also known as the Jewel Bug Nebula or the Magic Carpet Nebula, is a very young and dense planetary nebula located around from Earth in the constellation Cygnus. Discovered in 1878 by Édouard Stephan using the reflector at Marseille Observatory, it is one of the smallest planetary nebulae and by far the most extensively studied.
Observation
NGC 7027 is one of the visually brightest planetary nebulae. In a telescope with a 6" aperture at a magnification of around 50× it appears as a relatively bright bluish star. It is best viewed with the highest magnification possible.
In 1977 at Yerkes Observatory, a small Schmidt-Cassegrain telescope was used to derive an accurate optical position for the planetary nebula NGC 7027 to allow comparison between photographs and radio maps of the object.
It has been photographed multiple times by the Hubble Space Telescope since its launch in 1990. Prior to these observations, NGC 7027 was thought to be a protoplanetary nebula with the central star too cool to ionize any of the gas, but it is now known to be a planetary nebula in the earliest stage of its development.
Overview
NGC 7027 is unusually small, measuring only 0.2 by 0.1 light-years, whereas the typical size for a planetary nebula is 1 light-year. It is fairly young, at about 600 years old. It has a very complex shape, consisting of an elliptical region of ionized gas and an equatorial belt within a massive neutral cloud. The inner structure is surrounded by a translucent shroud of gas and dust. The nebula is shaped like a prolate ellipsoidal shell and contains a photodissociation region shaped like a "clover leaf". The inner shell is also punctured by several shocks and X-ray jets, leading to the "spike"-like structures. NGC 7027 is expanding at . The central regions of NGC 7027 have been found to emit X-rays, indicating very high temperatures. Surrounding the ellipsoidal nebula are a series of faint, blue concentric shells.
The expanding halo of NGC 7027 has a mass of about three times the mass of the Sun, and is about 100 times more massive than the ionized central region. This mass loss in NGC 7027 provided important evidence that stars a few times more massive than the Sun can avoid being destroyed in supernova explosions.
The nebula is rich in carbon, and is a very interesting object for the study of carbon chemistry in dense molecular material exposed to strong ultraviolet radiation. The spectrum of NGC 7027 contains fewer spectral lines from neutral molecules than is usual for planetary nebulae. This is due to the destruction of neutral molecules by intense UV radiation. The nebula contains ions of extremely high ionization potential. The helium hydride ion, thought to be the earliest molecule to have been formed in the Universe (about 100,000 years after the Big Bang), was detected in 2019 for the first time in space in NGC 7027. There is also evidence for the presence of nanodiamond in NGC 7027.
Central star
NGC 7027 has a rich and highly ionized spectrum caused by its hot central star. The progenitor star of NGC 7027 is believed to have been about 3 to 4 times the mass of the Sun before the nebula was formed. It is possible that the central white dwarf of NGC 7027 has an accretion disk that acts as a source of high temperatures. The white dwarf is believed to have a mass approximately 0.7 times the mass of the Sun and is radiating at 7,700 times the Sun's luminosity. NGC 7027 is currently in a short phase of planetary nebula evolution in which molecules in its envelope are being dissociated into their component atoms, and the atoms are being ionized.
The central star is suspected to be a binary system with the secondary being undetected. Although the details of NGC 7027's formation are unclear, it is hypothesized that interactions with the secondary star produced the complex structure of the planetary nebula, including the jets and resulting spikes.
| Physical sciences | Notable nebulae | Astronomy |
1338096 | https://en.wikipedia.org/wiki/Tilth | Tilth | Tilth is a physical condition of soil, especially in relation to its suitability for planting or growing a crop. Factors that determine tilth include the formation and stability of aggregated soil particles, moisture content, degree of aeration, soil biota, rate of water infiltration and drainage. Tilth can change rapidly, depending on environmental factors, such as changes in moisture, tillage and soil amendments. The objective of tillage (mechanical manipulation of the soil) is to improve tilth, thereby increasing crop production; in the long term, however, conventional tillage, especially plowing, often has the opposite effect, causing the soil carbon sponge to oxidize, break down and become compacted.
Soil with good tilth is spongy with large pore spaces for air infiltration and water movement. Roots grow only where the soil tilth allows for adequate levels of soil oxygen. Such soil also holds a reasonable supply of water and nutrients.
Tillage, organic matter amendments, fertilization and irrigation can each improve tilth, but when used excessively, can have the opposite effect. Crop rotation and cover crops can rebuild the soil carbon sponge and positively affect tilth. A combined approach can produce the greatest improvement.
Aggregation
Good tilth shares a balanced relation between soil-aggregate tensile strength and friability, in which it has a stable mixture of aggregate soil particles that can be readily broken up by shallow, non-abrasive tilling. A high tensile strength will result in large cemented clods of compacted soil with low friability. Proper management of agricultural soils can positively affect soil aggregation and improve tilth quality.
Aggregation is positively associated with tilth. With finer-textured soils, aggregates may in turn be made up of smaller aggregates. Aggregation implies substantial pores between individual aggregates.
Aggregation is important in the subsoil, the layer below tillage. Such aggregates involve larger (2- to 6-inch) blocks of soil that are more angular and not as distinctive. These aggregates are less affected by biological activity than the tillage layer. Subsurface aggregates are important for root growth deep into the profile. Deep roots allow greater access to moisture, which helps in drought periods. Subsoil aggregates can also be compacted, mainly by heavy equipment on wet soil. Another significant source of subsoil compaction is the practice of plowing with tractor wheels in the open furrow.
Pore size
Soil that is well aggregated has a range of pore sizes. Each pore size plays a role in soil's physical functioning. Large pores drain rapidly and are needed for good air exchange during wet periods, preventing oxygen deficiency that can drown plants and increase pest problems. Oxygen-deficient wet soils increase denitrification – conversion of nitrogen to gaseous forms. In degraded soil, large pores are compressed into small ones.
Small pores are critical for water retention and help a crop endure dry periods with minimal yield loss.
Management
Soil tilth is naturally maintained by the interaction of plant roots with the soil biota.
Short lived tilth can be obtained through mechanical and biological manipulation.
Tillage
In 2021, the globally tilled soil volume was estimated at 1840 km3/yr. This value exceeds by two orders of magnitude the global total of all engineering earthworks. For comparison globally, the natural process of soil bioturbation by plant roots and earthworms, was estimated at 960 km3/yr.
Mechanical soil cultivation practices, including primary tillage (mold-board or chisel plowing) followed by secondary tillage (disking, harrowing, etc.), break up and aerate soil. Mechanical traffic and intensive tilling methods have a negative impact on soil aggregates, friability, soil porosity, and soil-bulk density. When soils become degraded and compacted, such tillage practices are often deemed necessary. The tilth created by tillage, however, tends to be unstable, because the aggregation is obtained through the physical manipulation of the soil, which is short lived, especially after years of intensive tillage. The compaction of soil aggregates can also decrease soil biota due to the low levels of oxygen in the top-soil. The resulting high soil-bulk density results in lower water infiltration from rainfall or conventional irrigation (surface, sprinkler, center-pivot); in turn, the series of processes will naturally erode and dissolve small soil particles and organic matter. The consequences from these processes cyclically require more tilling and intervention, thus tillage practices have the capability to disrupt biological mechanisms that stabilize soil structure, the soil carbon sponge and tilth quality.
Biological
The preferred scenario for good tilth is as the result of natural soil-building processes, provided by the activity of plant roots, microorganisms, earthworms and other beneficial organisms. Such stable aggregates break apart during tillage/planting and readily provide good tilth. Soil biota and organic matter work in unison to bind soil aggregates and establish a natural soil stability – a soil carbon sponge. Plant root exudates feed bacteria that emit extracellular polysaccharides (EPS), and feed the growth of fungal hyphae, to form a soil carbon sponge with the dispersed clay particles. These active tilth-forming processes contribute to the formation and stabilization of soil structure. The resulting soil structure reduces tensile strength and soil-bulk density while still forming soil aggregates through their abiotic/biotic binding mechanisms that resist breakdown during water saturation. The fungal hyphae networks can establish a role of enmeshment with EPS and rhizodeposition, thus improving aggregate stability. However, these organic materials are themselves subject to biological degradation, requiring active amendments with organic material, and minimal mechanical tillage. Tilth quality is heavily dependent on these naturally binding processes between biotic microorganisms and abiotic soil particles, as well as the necessary input of organic matter. All constituents in this naturally binding network must be supplied or managed in agriculture to ensure the sustainability of their presence through growing seasons.
Rotation
Crop rotation can help restore tilth in compacted soils. Two processes contribute to this gain. First, accelerated organic matter decomposition from tillage ends under the sod crop. Another way to achieve this is via no-till farming. Second, grass and legume sods develop extensive root systems that continually grow and die off. The dead roots supply a source of active organic matter, which feeds soil organisms that create aggregation – the soil carbon sponge. Beneficial organisms need continual supplies of organic matter to sustain themselves and they deposit the digested materials on soil aggregates and thereby stabilize them. Also, the living roots and symbiotic microorganisms (for example, mycorrhizal fungi) can exude organic materials that nourish soil organisms and help with aggregation. Grass and legume sod crops therefore deposit more organic matter in the soil than most other crops.
Some annual rotation crops, such as buckwheat, also have dense, fibrous, root systems and can improve tilth. Crop mixtures with different rooting systems can be beneficial. For example, red clover seeded into winter wheat provides additional roots and a more protein-rich soil organic matter.
Other rotation crops are more valuable for improving subsoils. Perennial crops, such as alfalfa, have strong, deep, penetrating tap roots that can push through hard layers, especially during wet periods when the soil is soft. These deep roots establish pathways for water and future plant roots, and produce soil organic matter.
Crops rotation can extend the period of active growth compared to conventional row crops, leaving more organic material behind. For example, in a corn–soybean rotation, active growth occurs 32% of the time, while a dry bean–winter wheat–corn rotation is active 72% of the time. Crops such as rye, wheat, oat, barley, pea and cool-season grasses grow actively in the late fall and early spring when other crops are inactive. They are beneficial both as rotation and cover crops, although intensive tillage can negate their effects.
Soil types
The soil management practices required to maintain soil tilth are a function of the type of soil. Sandy and gravelly soils are naturally deficient in small pores and are therefore drought prone, whereas loams and clays can retain and thus supply crops with more water.
Coarse-textured, sandy soils
Sandy soil has lower capacity to hold water and nutrients. Water is applied more frequently in smaller amounts to avoid it leaching and carrying nutrients below the root zone. Routine application of organic matter increases sandy soil's ability to hold water and nutrients by 10 times or more.
Fine-textured, clay soils
Clay soils lack large pores, restricting both water and air movement. During irrigation or rain events, the limited large pore space in fine-textured soils quickly fills with water, reducing soil oxygen levels. In addition to routine application of organic matter, microorganisms and earthworms perform a crucial assist to soil tilth. As microorganisms decompose the organic matter, soil particles bind together into larger aggregates, increasing large pore space. Clay soils are more subject to soil compaction, which reduces large pore spaces.
Gravelly and decomposed granite soils
Such soils natively have little tilth, especially once they have been disturbed. Adding organic matter up to 25% by volume can help compensate. For example, if tilling to a depth of eight inches, add two inches of organic materials.
| Physical sciences | Soil science | Earth science |
1338741 | https://en.wikipedia.org/wiki/Peculiar%20galaxy | Peculiar galaxy | A peculiar galaxy is a galaxy of unusual size, shape, or composition. Between five and ten percent of known galaxies are categorized as peculiar. Astronomers have identified two types of peculiar galaxies: interacting galaxies and active galactic nuclei (AGN).
When two galaxies come close to each other, their mutual gravitational forces can cause them to acquire highly irregular shapes. The terms 'peculiar galaxy' and 'interacting galaxy' have now become synonymous because the majority of peculiar galaxies attribute their forms to such gravitational forces.
Formation
Scientists hypothesize that many peculiar galaxies are formed by the collision of two or more galaxies. As such, peculiar galaxies tend to host more active galactic nuclei than normal galaxies, indicating that they contain supermassive black holes. Many peculiar galaxies experience starbursts, or episodes of rapid star formation, due to the galaxies merging. The periods of elevated star formation and the luminosity resulting from active galactic nuclei cause peculiar galaxies to be slightly bluer in color than other galaxies.
Studying peculiar galaxies can offer insights on other types of galaxies by providing useful information on galactic formation and evolution. Arp mapped peculiar galaxies in his 1966 Atlas of Peculiar Galaxies. Arp states that "the peculiarities of the galaxies pictured in this Atlas represent perturbations, deformations, and interactions which should enable us to analyze the nature of the real galaxies which we observe and which are too remote to experiment on directly."
Notation
Peculiar galaxies are notated by an additional "p" or "pec" (depending on the exact convention) after the Hubble type of the galaxy.
| Physical sciences | Galaxy classification | Astronomy |
25516868 | https://en.wikipedia.org/wiki/Ochetellus%20glaber | Ochetellus glaber | Ochetellus glaber (also known as the black household ant) is a species of ant native to Australia. A member of the genus Ochetellus in the subfamily Dolichoderinae, it was described by Austrian entomologist Gustav Mayr in 1862. Aside from Australia, O. glaber has been introduced to a number of countries, including China, India, Japan, New Zealand, the Philippines and the United States, where it has established itself in Hawaii and Florida. It has been found on Lord Howe Island, New Caledonia, Norfolk Island, Réunion, New Zealand, and the Solomon Islands. Compared with other ants, O. glaber is a small species, with workers measuring . Males are the smallest at , while the queens measure . The ant's colour ranges from brown to black.
Described as an arboreal nesting species, O. glaber lives in open or savannah woodland areas, nesting under stones, old dry logs, rotten wood, and in hollow trees and plant stems. Nests can also be constructed in buildings and structures, specifically in pavings, ceilings and walls. It is both diurnal and nocturnal, forming long trails from trees in search of food such as honeydew and insects. It has developed some associations with certain flowers and also tends to associate with some insects, such as mealybugs and aphids. During the nuptial flight, queens mate with either one or multiple males; males only mate with a single queen. Sometimes, a subset of a colony may leave the main colony for an alternative nest site as an act of dispersal. O. glaber often invades human homes to feed on household foods, and is considered a household pest. It has been intercepted numerous times in the United States, where it has the potential to disrupt the biological control of certain pests and cause long-term ecological impacts in areas where it is not native.
Taxonomy
Ochetellus glaber was originally described as Hypoclinea glabra by Austrian entomologist Gustav Mayr in his "Myrmecologische studien" in 1862. The ant was described from syntype workers and males Mayr collected from Sydney, Australia, now preserved in the Natural History Museum, London. Its placement in Hypoclinea was relatively short, as Mayr transferred it to the genus Iridomyrmex as Iridomyrmex glaber in 1865. Its placement in Iridomyrmex was accepted for more than a century, until entomologist Steve Shattuck revised the genus in 1992 and transferred the ant into a new genus, Ochetellus. He also designated O. glaber as the type species of the genus. In 2011, evidence emerged that O. glaber represents a species complex, indicating that the current taxa may need to be split.
O. glaber has two synonyms, Iridomyrmex itoi and Iridomyrmex itoi abbotti. I. itoi was described as a separate species, while I. itoi abbotti was recognised as a subspecies. The subspecies was short-lived, and was synonymised with I. itoi in 1910, and in the 1950s I. itoi was first noted to be similar to O. glaber. Despite the similarities, I. itoi remained a valid species until a 1995 publication confirmed its synonymy with O. glaber. It is commonly known as the black house ant or the tramp ant.
Description
O. glaber is small, with workers measuring . The antenna has 12 segments, the scapes of which are half as long as the head. Its antennal sockets and posterior clypeal margin are separated from one another by a small distance, perhaps less than the minimum width of the antennal scapes. Eyes range in size, being either medium to large, with more than six facets (lenses that make up the compound eye of an insect). The dorsum of the mesosoma has distinct metanotal grooves and lack erect hairs. The propodeum has a distinct protrusion which causes the slope to be strongly concave. The ant's waist has only one segment. The petiole (the narrow waist) is upright and is not flattened. The gaster has a ventral slit. Constriction between the third and fourth abdominal segment is not visible. The ant's colour ranges from brown to black. Males are smaller than the workers, measuring . The body is brown in colour, but the back of the body is brownish-black, and the mandibles, legs and antennae are yellow. The head and thorax are noticeably wrinkled. O. glaber queens are larger than the males and workers, measuring .
A young larva is . Compared to older larvae, young larva bodies are stouter and the outlines are straighter. Thirteen differentiated somites are present. Spinules (small spines or thorns) are more prominent at the posterior end. The length of the body hair is extremely short, measuring 0.002–0.015mm long. Mature larvae are larger, measuring . The body is short and stout. The integument (a tough outer protective layer of an organism) are spinulose, bearing small spines. These spinules are in short transverse rows both ventrally and posteriorly. Body hair is significant but few hairs are found on the head. Mandibles contain a large apical tooth. The maxillary and labial palps (organs which aid sensory function in eating) have three sensilla (a sensory organ protruding from the cuticle). Unlike other Dolichoderines, the larvae are yellow, not white.
The karyotype of O. glaber has been described. It has eight metacentric chromosomes, four submetacentric to acrocentric, and two submetacentric chromosomes.
Subspecies
Three subspecies of O. glaber are recognised: O. glaber clarithorax, O. glaber consimilis and O. glaber sommeri.
Distribution and habitat
O. glaber inhabits many areas in Oceania. In Australia, its range extends from coastal Queensland and New South Wales to south-west Western Australia. O. glaber was introduced to New Zealand around 1927 and had become well-established by the 1940s. While the ant has largely remained in Auckland and some suburbs, New Zealand authorities have intercepted specimens elsewhere several times, and it may spread to other New Zealand cities. It is regarded as a potential pest, though not a major household one. Additionally, specimens of O. glaber have been collected on Lord Howe Island, New Caledonia, Norfolk Island, Solomon Islands and the former New Hebrides.
Other locations where O. glaber has been found include Réunion, India and the Philippines. In India, it has been collected from the states of Haryana, Himachal Pradesh, Karnataka, Maharashtra and Uttarakhand. It has also been found in China and Macao, as well as Japan and Sri Lanka. In the United States, O. glaber was first recorded in Hawaii in 1977, where it originated from Australia and Japan. In Hawaii, it is presently found on Kauai, Maui, Oahu, and the Island of Hawaii. It was also found in Florida, where collected specimens were found in a queen palm (Syagrus romanzoffiana) stump. It is abundant yet localised species in Orange County, being found in dead wood or in tussocks of marsh grasses.
O. glaber is an arboreal nesting species. It lives in open areas or savannah woodland, nesting under stones or old dry logs, or else in hollow trees, plant stems or rotten wood. It is also often found in gardens, where it may be conspicuous. O. glaber has also been found in mountain forests, wet forests, in pastures, garden flower tubs and dried palmetto frond. It is found at altitudes of between above sea level. In buildings and structures, O. glaber nests in crevices and cavities such as rockeries, paving and in brickwork. It also nests in ceilings, walls, and subfloor areas.
Behaviour and ecology
O. glaber is omnivorous, forming long trails on tree trunks to seek sweet substances such as honeydew and to hunt insects. The ant is both active during the day and night. Activity increases during the night or on overcast days, peaking during early mornings and late evening to early night. Nocturnal activity varies but is either minimal or non-existent. O. glaber consumes carcasses of dead birds, sea turtles, parrot fish, fruit fly pupae and diamondback moth larvae. It also has a preference for fat, grease, plants and seeds. Galleries of the Formosan subterranean termite (Coptotermes formosanus) are often invaded by O. glaber, although mortality rates are substantially higher when big-headed ants (Pheidole megacephala) invade termite galleries.
The ant has developed associations with a range of organisms. Foraging workers often visit flowers for nectar, chiefly those of Pisonia, but also of Canavalia, Commicarpus, Ipomoea, Melanthera, Plumbago and Scaevola. O. glaber also associates with some insects, including the pineapple mealybug (Dysmicoccus brevipes) and aphids, which they import and tend with other bugs on domestic pot plants. O. glaber associates indirectly with Ananusia australis, an encyrtid parasitoid wasp.
During the nuptial flight, a queen may mate with multiple males while the males will only mate with a single queen, making them monogynous. Sometimes, colonies proliferate by "budding" (also called "satelliting" or "fractionating") whereby a subset of the colony, including queens, workers and brood (eggs, larvae and pupae) leave the main colony for an alternative nest site.
Relationship with humans
O. glaber is viewed as a pest. Although it does not sting, it bites and, when crushed, produces a strong odour. It enters human homes to gather food, tracking across ceilings, beams and joists and drops ant debris onto surfaces below. The ant can also quickly spread because of its high reproduction and dispersal potential.
O. glaber is more potentially dangerous in some places than in others. The United States, and especially the state of California, prohibits the ant from entering lest it should disrupt the country's ecosystems. Its relationship with honeydew-producing insects and consumption of parasitoids could disrupt the biological control of certain pests, risking long-term ecological damage that could reduce biodiversity, disrupt natural communities, or change ecosystem processes. Therefore, the California Department of Food and Agriculture have intercepted the ant in nursery stock and fresh plants from Hawaii. California is especially vulnerable to O. glaber infestation because the state's climate resembles that of those regions where the ant already lives. Nonetheless, O. glaber is unlikely to lower crop yields, increase farming costs, degrade water supplies, or likely disrupt Californian agricultural commodity markets.
In New Zealand, the ant is found only in urban gardens and some homes.
| Biology and health sciences | Hymenoptera | Animals |
21144218 | https://en.wikipedia.org/wiki/Electron%20mass | Electron mass | In particle physics, the electron mass (symbol: ) is the mass of a stationary electron, also known as the invariant mass of the electron. It is one of the fundamental constants of physics. It has a value of about or about , which has an energy-equivalent of about or about .
Terminology
The term "rest mass" is sometimes used because in special relativity the mass of an object can be said to increase in a frame of reference that is moving relative to that object (or if the object is moving in a given frame of reference). Most practical measurements are carried out on moving electrons. If the electron is moving at a relativistic velocity, any measurement must use the correct expression for mass. Such correction becomes substantial for electrons accelerated by voltages of over .
For example, the relativistic expression for the total energy, , of an electron moving at speed is
where
is the speed of light;
is the Lorentz factor,
is the "rest mass", or more simply just the "mass" of the electron.
This quantity is frame invariant and velocity independent.
Determination
Since the electron mass determines a number of observed effects in atomic physics, there are potentially many ways to determine its mass from an experiment, if the values of other physical constants are already considered known.
Historically, the mass of the electron was determined directly from combining two measurements. The mass-to-charge ratio of the electron was first estimated by Arthur Schuster in 1890 by measuring the deflection of "cathode rays" due to a known magnetic field in a cathode ray tube. Seven years later J. J. Thomson showed that cathode rays consist of streams of particles, to be called electrons, and made more precise measurements of their mass-to-charge ratio again using a cathode ray tube.
The second measurement was of the charge of the electron. This was determined with a precision of better than 1% by Robert A. Millikan in his oil drop experiment in 1909. Together with the mass-to-charge ratio, the electron mass was determined with reasonable precision. The value of mass that was found for the electron was initially met with surprise by physicists, since it was so small (less than 0.1%) compared to the known mass of a hydrogen atom.
The electron rest mass can be calculated from the Rydberg constant and the fine-structure constant obtained through spectroscopic measurements. Using the definition of the Rydberg constant:
thus
where is the speed of light and is the Planck constant. The relative uncertainty, 5 in the 2006 CODATA recommended value, is due entirely to the uncertainty in the value of the Planck constant. With the re-definition of kilogram in 2019, there is no uncertainty by definition left in Planck constant anymore.
The electron relative atomic mass can be measured directly in a Penning trap. It can also be inferred from the spectra of antiprotonic helium atoms (helium atoms where one of the electrons has been replaced by an antiproton) or from measurements of the electron g-factor in the hydrogenic ions 12C5+ or 16O7+.
The electron relative atomic mass is an adjusted parameter in the CODATA set of fundamental physical constants, while the electron rest mass in kilograms is calculated from the values of the Planck constant, the fine-structure constant and the Rydberg constant, as detailed above.
Relationship to other physical constants
The electron mass was used to calculate the Avogadro constant before its value was fixed as a defining constant in the 2019 revision of the SI:
Hence it is also related to the atomic mass constant :
where
is the molar mass constant (defined in SI);
is a directly measured quantity, the relative atomic mass of the electron.
is defined in terms of , and not the other way round, and so the name "electron mass in atomic mass units" for involves a circular definition (at least in terms of practical measurements).
The electron relative atomic mass also enters into the calculation of all other relative atomic masses. By convention, relative atomic masses are quoted for neutral atoms, but the actual measurements are made on positive ions, either in a mass spectrometer or a Penning trap. Hence the mass of the electrons must be added back on to the measured values before tabulation. A correction must also be made for the mass equivalent of the binding energy . Taking the simplest case of complete ionization of all electrons, for a nuclide X of atomic number ,
As relative atomic masses are measured as ratios of masses, the corrections must be applied to both ions: the uncertainties in the corrections are negligible, as illustrated below for hydrogen 1 and oxygen 16.
The principle can be shown by the determination of the electron relative atomic mass by Farnham et al. at the University of Washington (1995). It involves the measurement of the frequencies of the cyclotron radiation emitted by electrons and by ions in a Penning trap. The ratio of the two frequencies is equal to six times the inverse ratio of the masses of the two particles (the heavier the particle, the lower the frequency of the cyclotron radiation; the higher the charge on the particle, the higher the frequency):
As the relative atomic mass of ions is very nearly 12, the ratio of frequencies can be used to calculate a first approximation to Ar(e), . This approximate value is then used to calculate a first approximation to Ar(12C6+), knowing that (from the sum of the six ionization energies of carbon) is : . This value is then used to calculate a new approximation to Ar(e), and the process repeated until the values no longer vary (given the relative uncertainty of the measurement, 2.1): this happens by the fourth cycle of iterations for these results, giving for these data.
| Physical sciences | Mass and weight | Basics and measurement |
2719753 | https://en.wikipedia.org/wiki/Scattered%20disc | Scattered disc | The scattered disc (or scattered disk) is a distant circumstellar disc in the Solar System that is sparsely populated by icy small Solar System bodies, which are a subset of the broader family of trans-Neptunian objects. The scattered-disc objects (SDOs) have orbital eccentricities ranging as high as 0.8, inclinations as high as 40°, and perihelia greater than . These extreme orbits are thought to be the result of gravitational "scattering" by the gas giants, and the objects continue to be subject to perturbation by the planet Neptune.
Although the closest scattered-disc objects approach the Sun at about 30–35 AU, their orbits can extend well beyond 100 AU. This makes scattered disc objects among the coldest and most distant objects in the Solar System. The innermost portion of the scattered disc overlaps with a torus-shaped region of orbiting objects traditionally called the Kuiper belt, but its outer limits reach much farther away from the Sun and farther above and below the ecliptic than the Kuiper belt proper.
Because of its unstable nature, astronomers now consider the scattered disc to be the place of origin for most periodic comets in the Solar System, with the centaurs, a population of icy bodies between Jupiter and Neptune, being the intermediate stage in an object's migration from the disc to the inner Solar System. Eventually, perturbations from the giant planets send such objects towards the Sun, transforming them into periodic comets. Many objects of the proposed Oort cloud are also thought to have originated in the scattered disc. Detached objects are not sharply distinct from scattered disc objects, and some such as Sedna have sometimes been considered to be included in this group.
Discovery
Traditionally, devices like a blink comparator were used in astronomy to detect objects in the Solar System, because these objects would move between two exposures—this involved time-consuming steps like exposing and developing photographic plates or films, and people then using a blink comparator to manually detect prospective objects. During the 1980s, the use of CCD-based cameras in telescopes made it possible to directly produce electronic images that could then be readily digitized and transferred to digital images. Because the CCD captured more light than film (about 90% versus 10% of incoming light) and the blinking could now be done at an adjustable computer screen, the surveys allowed for higher throughput. A flood of new discoveries was the result: over a thousand trans-Neptunian objects were detected between 1992 and 2006.
The first scattered-disc object (SDO) to be recognised as such was , originally identified in 1996 by astronomers based at Mauna Kea in Hawaii. Three more were identified by the same survey in 1999: , , and . The first object presently classified as an SDO to be discovered was , found in 1995 by Spacewatch.
As of 2011, over 200 SDOs have been identified, including Gǃkúnǁʼhòmdímà (discovered by Schwamb, Brown, and Rabinowitz), Gonggong (Schwamb, Brown, and Rabinowitz) (NEAT), Eris (Brown, Trujillo, and Rabinowitz), Sedna (Brown, Trujillo, and Rabinowitz), and 474640 Alicanto (Deep Ecliptic Survey). Although the numbers of objects in the Kuiper belt and the scattered disc are hypothesized to be roughly equal, observational bias due to their greater distance means that far fewer SDOs have been observed to date.
Subdivisions of trans-Neptunian space
Known trans-Neptunian objects are often divided into two subpopulations: the Kuiper belt and the scattered disc. A third reservoir of trans-Neptunian objects, the Oort cloud, has been hypothesized, although no confirmed direct observations of the Oort cloud have been made. Some researchers further suggest a transitional space between the scattered disc and the inner Oort cloud, populated with "detached objects".
Scattered disc versus Kuiper belt
The Kuiper belt is a relatively thick torus (or "doughnut") of space, extending from about 30 to 50 AU comprising two main populations of Kuiper belt objects (KBOs): the classical Kuiper-belt objects (or "cubewanos"), which lie in orbits untouched by Neptune, and the resonant Kuiper-belt objects, those which Neptune has locked into a precise orbital ratio such as 2:3 (the object goes around twice for every three Neptune orbits) and 1:2 (the object goes around once for every two Neptune orbits). These ratios, called orbital resonances, allow KBOs to persist in regions which Neptune's gravitational influence would otherwise have cleared out over the age of the Solar System, since the objects are never close enough to Neptune to be scattered by its gravity. Those in 2:3 resonances are known as "plutinos", because Pluto is the largest member of their group, whereas those in 1:2 resonances are known as "twotinos".
In contrast to the Kuiper belt, the scattered-disc population can be disturbed by Neptune. Scattered-disc objects come within gravitational range of Neptune at their closest approaches (~30 AU) but their farthest distances reach many times that. Ongoing research suggests that the centaurs, a class of icy planetoids that orbit between Jupiter and Neptune, may simply be SDOs thrown into the inner reaches of the Solar System by Neptune, making them "cis-Neptunian" rather than trans-Neptunian scattered objects. Some objects, like (29981) 1999 TD10, blur the distinction and the Minor Planet Center (MPC), which officially catalogues all trans-Neptunian objects, now lists centaurs and SDOs together.
The MPC, however, makes a clear distinction between the Kuiper belt and the scattered disc, separating those objects in stable orbits (the Kuiper belt) from those in scattered orbits (the scattered disc and the centaurs). However, the difference between the Kuiper belt and the scattered disc is not clear-cut, and many astronomers see the scattered disc not as a separate population but as an outward region of the Kuiper belt. Another term used is "scattered Kuiper-belt object" (or SKBO) for bodies of the scattered disc.
Morbidelli and Brown propose that the difference between objects in the Kuiper belt and scattered-disc objects is that the latter bodies "are transported in semi-major axis by close and distant encounters with Neptune," but the former experienced no such close encounters. This delineation is inadequate (as they note) over the age of the Solar System, since bodies "trapped in resonances" could "pass from a scattering phase to a non-scattering phase (and vice versa) numerous times." That is, trans-Neptunian objects could travel back and forth between the Kuiper belt and the scattered disc over time. Therefore, they chose instead to define the regions, rather than the objects, defining the scattered disc as "the region of orbital space that can be visited by bodies that have encountered Neptune" within the radius of a Hill sphere, and the Kuiper belt as its "complement ... in the a > 30 AU region"; the region of the Solar System populated by objects with semi-major axes greater than 30 AU.
Detached objects
The Minor Planet Center classifies the trans-Neptunian object 90377 Sedna as a scattered-disc object. Its discoverer Michael E. Brown has suggested instead that it should be considered an inner Oort-cloud object rather than a member of the scattered disc, because, with a perihelion distance of 76 AU, it is too remote to be affected by the gravitational attraction of the outer planets. Under this definition, an object with a perihelion greater than 40 AU could be classified as outside the scattered disc.
Sedna is not the only such object: (discovered before Sedna) and 474640 Alicanto have a perihelion too far away from Neptune to be influenced by it. This led to a discussion among astronomers about a new minor planet set, called the extended scattered disc (E-SDO). may also be an inner Oort-cloud object or (more likely) a transitional object between the scattered disc and the inner Oort cloud. More recently, these objects have been referred to as "detached", or distant detached objects (DDO).
There are no clear boundaries between the scattered and detached regions. Gomes et al. define SDOs as having "highly eccentric orbits, perihelia beyond Neptune, and semi-major axes beyond the 1:2 resonance." By this definition, all distant detached objects are SDOs. Since detached objects' orbits cannot be produced by Neptune scattering, alternative scattering mechanisms have been put forward, including a passing star or a distant, planet-sized object. Alternatively, it has been suggested that these objects have been captured from a passing star.
A scheme introduced by a 2005 report from the Deep Ecliptic Survey by J. L. Elliott et al. distinguishes between two categories: scattered-near (i.e. typical SDOs) and scattered-extended (i.e. detached objects). Scattered-near objects are those whose orbits are non-resonant, non-planetary-orbit-crossing and have a Tisserand parameter (relative to Neptune) less than 3. Scattered-extended objects have a Tisserand parameter (relative to Neptune) greater than 3 and have a time-averaged eccentricity greater than 0.2.
An alternative classification, introduced by B. J. Gladman, B. G. Marsden and C. Van Laerhoven in 2007, uses 10-million-year orbit integration instead of the Tisserand parameter. An object qualifies as an SDO if its orbit is not resonant, has a semi-major axis no greater than 2000 AU, and, during the integration, its semi-major axis shows an excursion of 1.5 AU or more. Gladman et al. suggest the term scattering disk object to emphasize this present mobility. If the object is not an SDO as per the above definition, but the eccentricity of its orbit is greater than 0.240, it is classified as a detached TNO. (Objects with smaller eccentricity are considered classical.) In this scheme, the disc extends from the orbit of Neptune to 2000 AU, the region referred to as the inner Oort cloud.
Orbits
The scattered disc is a very dynamic environment. Because they are still capable of being perturbed by Neptune, SDOs' orbits are always in danger of disruption; either of being sent outward to the Oort cloud or inward into the centaur population and ultimately the Jupiter family of comets. For this reason Gladman et al. prefer to refer to the region as the scattering disc, rather than scattered. Unlike Kuiper-belt objects (KBOs), the orbits of scattered-disc objects can be inclined as much as 40° from the ecliptic.
SDOs are typically characterized by orbits with medium and high eccentricities with a semi-major axis greater than 50 AU, but their perihelia bring them within influence of Neptune. Having a perihelion of roughly 30 AU is one of the defining characteristics of scattered objects, as it allows Neptune to exert its gravitational influence.
The classical objects (cubewanos) are very different from the scattered objects: more than 30% of all cubewanos are on low-inclination, near-circular orbits whose eccentricities peak at 0.25. Classical objects possess eccentricities ranging from 0.2 to 0.8. Though the inclinations of scattered objects are similar to the more extreme KBOs, very few scattered objects have orbits as close to the ecliptic as much of the KBO population.
Although motions in the scattered disc are random, they do tend to follow similar directions, which means that SDOs can become trapped in temporary resonances with Neptune. Examples of possible resonant orbits within the scattered disc include 1:3, 2:7, 3:11, 5:22 and 4:79.
Formation
The scattered disc is still poorly understood: no model of the formation of the Kuiper belt and the scattered disc has yet been proposed that explains all their observed properties.
According to contemporary models, the scattered disc formed when Kuiper belt objects (KBOs) were "scattered" into eccentric and inclined orbits by gravitational interaction with Neptune and the other outer planets. The amount of time for this process to occur remains uncertain. One hypothesis estimates a period equal to the entire age of the Solar System; a second posits that the scattering took place relatively quickly, during Neptune's early migration epoch.
Models for a continuous formation throughout the age of the Solar System illustrate that at weak resonances within the Kuiper belt (such as 5:7 or 8:1), or at the boundaries of stronger resonances, objects can develop weak orbital instabilities over millions of years. The 4:7 resonance in particular has large instability. KBOs can also be shifted into unstable orbits by close passage of massive objects, or through collisions. Over time, the scattered disc would gradually form from these isolated events.
Computer simulations have also suggested a more rapid and earlier formation for the scattered disc. Modern theories indicate that neither Uranus nor Neptune could have formed in situ beyond Saturn, as too little primordial matter existed at that range to produce objects of such high mass. Instead, these planets, and Saturn, may have formed closer to Jupiter, but were flung outwards during the early evolution of the Solar System, perhaps through exchanges of angular momentum with scattered objects. Once the orbits of Jupiter and Saturn shifted to a 2:1 resonance (two Jupiter orbits for each orbit of Saturn), their combined gravitational pull disrupted the orbits of Uranus and Neptune, sending Neptune into the temporary "chaos" of the proto-Kuiper belt. As Neptune traveled outward, it scattered many trans-Neptunian objects into higher and more eccentric orbits. This model states that 90% or more of the objects in the scattered disc may have been "promoted into these eccentric orbits by Neptune's resonances during the migration epoch...[therefore] the scattered disc might not be so scattered."
Composition
Scattered objects, like other trans-Neptunian objects, have low densities and are composed largely of frozen volatiles such as water and methane. Spectral analysis of selected Kuiper belt and scattered objects has revealed signatures of similar compounds. Both Pluto and Eris, for instance, show signatures for methane.
Astronomers originally supposed that the entire trans-Neptunian population would show a similar red surface colour, as they were thought to have originated in the same region and subjected to the same physical processes. Specifically, SDOs were expected to have large amounts of surface methane, chemically altered into tholins by sunlight from the Sun. This would absorb blue light, creating a reddish hue. Most classical objects display this colour, but scattered objects do not; instead, they present a white or greyish appearance.
One explanation is the exposure of whiter subsurface layers by impacts; another is that the scattered objects' greater distance from the Sun creates a composition gradient, analogous to the composition gradient of the terrestrial and gas giant planets. Michael E. Brown, discoverer of the scattered object Eris, suggests that its paler colour could be because, at its current distance from the Sun, its atmosphere of methane is frozen over its entire surface, creating an inches-thick layer of bright white ice. Pluto, conversely, being closer to the Sun, would be warm enough that methane would freeze only onto cooler, high-albedo regions, leaving low-albedo tholin-covered regions bare of ice.
Comets
The Kuiper belt was initially thought to be the source of the Solar System's ecliptic comets. However, studies of the region since 1992 have shown that the orbits within the Kuiper belt are relatively stable, and that ecliptic comets originate from the scattered disc, where orbits are generally less stable.
Comets can loosely be divided into two categories: short-period and long-period—the latter being thought to originate in the Oort cloud. The two major categories of short-period comets are Jupiter-family comets (JFCs) and Halley-type comets. Halley-type comets, which are named after their prototype, Halley's Comet, are thought to have originated in the Oort cloud but to have been drawn into the inner Solar System by the gravity of the giant planets, whereas the JFCs are thought to have originated in the scattered disc. The centaurs are thought to be a dynamically intermediate stage between the scattered disc and the Jupiter family.
There are many differences between SDOs and JFCs, even though many of the Jupiter-family comets may have originated in the scattered disc. Although the centaurs share a reddish or neutral coloration with many SDOs, their nuclei are bluer, indicating a fundamental chemical or physical difference. One hypothesis is that comet nuclei are resurfaced as they approach the Sun by subsurface materials which subsequently bury the older material.
| Physical sciences | Solar System | Astronomy |
2720954 | https://en.wikipedia.org/wiki/Data%20analysis | Data analysis | Data analysis is the process of inspecting, cleansing, transforming, and modeling data with the goal of discovering useful information, informing conclusions, and supporting decision-making. Data analysis has multiple facets and approaches, encompassing diverse techniques under a variety of names, and is used in different business, science, and social science domains. In today's business world, data analysis plays a role in making decisions more scientific and helping businesses operate more effectively.
Data mining is a particular data analysis technique that focuses on statistical modeling and knowledge discovery for predictive rather than purely descriptive purposes, while business intelligence covers data analysis that relies heavily on aggregation, focusing mainly on business information. In statistical applications, data analysis can be divided into descriptive statistics, exploratory data analysis (EDA), and confirmatory data analysis (CDA). EDA focuses on discovering new features in the data while CDA focuses on confirming or falsifying existing hypotheses. Predictive analytics focuses on the application of statistical models for predictive forecasting or classification, while text analytics applies statistical, linguistic, and structural techniques to extract and classify information from textual sources, a species of unstructured data. All of the above are varieties of data analysis.
Data integration is a precursor to data analysis, and data analysis is closely linked to data visualization and data dissemination.
Data analysis process
Analysis refers to dividing a whole into its separate components for individual examination. Data analysis is a process for obtaining raw data, and subsequently converting it into information useful for decision-making by users. Data is collected and analyzed to answer questions, test hypotheses, or disprove theories.
Statistician John Tukey, defined data analysis in 1961, as:"Procedures for analyzing data, techniques for interpreting the results of such procedures, ways of planning the gathering of data to make its analysis easier, more precise or more accurate, and all the machinery and results of (mathematical) statistics which apply to analyzing data."There are several phases that can be distinguished, described below. The phases are iterative, in that feedback from later phases may result in additional work in earlier phases. The CRISP framework, used in data mining, has similar steps.
Data requirements
The data is necessary as inputs to the analysis, which is specified based upon the requirements of those directing the analytics (or customers, who will use the finished product of the analysis). The general type of entity upon which the data will be collected is referred to as an experimental unit (e.g., a person or population of people). Specific variables regarding a population (e.g., age and income) may be specified and obtained. Data may be numerical or categorical (i.e., a text label for numbers).
Data collection
Data is collected from a variety of sources. A list of data sources are available for study & research. The requirements may be communicated by analysts to custodians of the data; such as, Information Technology personnel within an organization. Data collection or data gathering is the process of gathering and measuring information on targeted variables in an established system, which then enables one to answer relevant questions and evaluate outcomes. The data may also be collected from sensors in the environment, including traffic cameras, satellites, recording devices, etc. It may also be obtained through interviews, downloads from online sources, or reading documentation.
Data processing
Data, when initially obtained, must be processed or organized for analysis. For instance, these may involve placing data into rows and columns in a table format (known as structured data) for further analysis, often through the use of spreadsheet(excel) or statistical software.
Data cleaning
Once processed and organized, the data may be incomplete, contain duplicates, or contain errors. The need for data cleaning will arise from problems in the way that the datum are entered and stored. Data cleaning is the process of preventing and correcting these errors. Common tasks include record matching, identifying inaccuracy of data, overall quality of existing data, deduplication, and column segmentation. Such data problems can also be identified through a variety of analytical techniques. For example; with financial information, the totals for particular variables may be compared against separately published numbers that are believed to be reliable. Unusual amounts, above or below predetermined thresholds, may also be reviewed. There are several types of data cleaning, that are dependent upon the type of data in the set; this could be phone numbers, email addresses, employers, or other values. Quantitative data methods for outlier detection, can be used to get rid of data that appears to have a higher likelihood of being input incorrectly. Textual data spell checkers can be used to lessen the amount of mistyped words. However, it is harder to tell if the words themselves are correct.
Exploratory data analysis
Once the datasets are cleaned, they can then be analyzed. Analysts may apply a variety of techniques, referred to as exploratory data analysis, to begin understanding the messages contained within the obtained data. The process of data exploration may result in additional data cleaning or additional requests for data; thus, the initialization of the iterative phases mentioned in the lead paragraph of this section. Descriptive statistics, such as, the average or median, can be generated to aid in understanding the data. Data visualization is also a technique used, in which the analyst is able to examine the data in a graphical format in order to obtain additional insights, regarding the messages within the data.
Modeling and algorithms
Mathematical formulas or models (also known as algorithms), may be applied to the data in order to identify relationships among the variables; for example, using correlation or causation. In general terms, models may be developed to evaluate a specific variable based on other variable(s) contained within the dataset, with some residual error depending on the implemented model's accuracy (e.g., Data = Model + Error).
Inferential statistics includes utilizing techniques that measure the relationships between particular variables. For example, regression analysis may be used to model whether a change in advertising (independent variable X), provides an explanation for the variation in sales (dependent variable Y). In mathematical terms, Y (sales) is a function of X (advertising). It may be described as (Y = aX + b + error), where the model is designed such that (a) and (b) minimize the error when the model predicts Y for a given range of values of X. Analysts may also attempt to build models that are descriptive of the data, in an aim to simplify analysis and communicate results.
Data product
A data product is a computer application that takes data inputs and generates outputs, feeding them back into the environment. It may be based on a model or algorithm. For instance, an application that analyzes data about customer purchase history, and uses the results to recommend other purchases the customer might enjoy.
Communication
Once data is analyzed, it may be reported in many formats to the users of the analysis to support their requirements. The users may have feedback, which results in additional analysis. As such, much of the analytical cycle is iterative.
When determining how to communicate the results, the analyst may consider implementing a variety of data visualization techniques to help communicate the message more clearly and efficiently to the audience. Data visualization uses information displays (graphics such as, tables and charts) to help communicate key messages contained in the data. Tables are a valuable tool by enabling the ability of a user to query and focus on specific numbers; while charts (e.g., bar charts or line charts), may help explain the quantitative messages contained in the data.
Quantitative messages
Stephen Few described eight types of quantitative messages that users may attempt to understand or communicate from a set of data and the associated graphs used to help communicate the message. Customers specifying requirements and analysts performing the data analysis may consider these messages during the course of the process.
Time-series: A single variable is captured over a period of time, such as the unemployment rate over a 10-year period. A line chart may be used to demonstrate the trend.
Ranking: Categorical subdivisions are ranked in ascending or descending order, such as a ranking of sales performance (the measure) by salespersons (the category, with each salesperson a categorical subdivision) during a single period. A bar chart may be used to show the comparison across the salespersons.
Part-to-whole: Categorical subdivisions are measured as a ratio to the whole (i.e., a percentage out of 100%). A pie chart or bar chart can show the comparison of ratios, such as the market share represented by competitors in a market.
Deviation: Categorical subdivisions are compared against a reference, such as a comparison of actual vs. budget expenses for several departments of a business for a given time period. A bar chart can show the comparison of the actual versus the reference amount.
Frequency distribution: Shows the number of observations of a particular variable for a given interval, such as the number of years in which the stock market return is between intervals such as 0–10%, 11–20%, etc. A histogram, a type of bar chart, may be used for this analysis.
Correlation: Comparison between observations represented by two variables (X,Y) to determine if they tend to move in the same or opposite directions. For example, plotting unemployment (X) and inflation (Y) for a sample of months. A scatter plot is typically used for this message.
Nominal comparison: Comparing categorical subdivisions in no particular order, such as the sales volume by product code. A bar chart may be used for this comparison.
Geographic or geospatial: Comparison of a variable across a map or layout, such as the unemployment rate by state or the number of persons on the various floors of a building. A cartogram is a typical graphic used.
Analyzing quantitative data
Author Jonathan Koomey has recommended a series of best practices for understanding quantitative data. These include:
Check raw data for anomalies prior to performing an analysis;
Re-perform important calculations, such as verifying columns of data that are formula driven;
Confirm main totals are the sum of subtotals;
Check relationships between numbers that should be related in a predictable way, such as ratios over time;
Normalize numbers to make comparisons easier, such as analyzing amounts per person or relative to GDP or as an index value relative to a base year;
Break problems into component parts by analyzing factors that led to the results, such as DuPont analysis of return on equity.
For the variables under examination, analysts typically obtain descriptive statistics for them, such as the mean (average), median, and standard deviation. They may also analyze the distribution of the key variables to see how the individual values cluster around the mean.
The consultants at McKinsey and Company named a technique for breaking a quantitative problem down into its component parts called the MECE principle. Each layer can be broken down into its components; each of the sub-components must be mutually exclusive of each other and collectively add up to the layer above them. The relationship is referred to as "Mutually Exclusive and Collectively Exhaustive" or MECE. For example, profit by definition can be broken down into total revenue and total cost. In turn, total revenue can be analyzed by its components, such as the revenue of divisions A, B, and C (which are mutually exclusive of each other) and should add to the total revenue (collectively exhaustive).
Analysts may use robust statistical measurements to solve certain analytical problems. Hypothesis testing is used when a particular hypothesis about the true state of affairs is made by the analyst and data is gathered to determine whether that state of affairs is true or false. For example, the hypothesis might be that "Unemployment has no effect on inflation", which relates to an economics concept called the Phillips Curve. Hypothesis testing involves considering the likelihood of Type I and type II errors, which relate to whether the data supports accepting or rejecting the hypothesis.
Regression analysis may be used when the analyst is trying to determine the extent to which independent variable X affects dependent variable Y (e.g., "To what extent do changes in the unemployment rate (X) affect the inflation rate (Y)?"). This is an attempt to model or fit an equation line or curve to the data, such that Y is a function of X.
Necessary condition analysis (NCA) may be used when the analyst is trying to determine the extent to which independent variable X allows variable Y (e.g., "To what extent is a certain unemployment rate (X) necessary for a certain inflation rate (Y)?"). Whereas (multiple) regression analysis uses additive logic where each X-variable can produce the outcome and the X's can compensate for each other (they are sufficient but not necessary), necessary condition analysis (NCA) uses necessity logic, where one or more X-variables allow the outcome to exist, but may not produce it (they are necessary but not sufficient). Each single necessary condition must be present and compensation is not possible.
Analytical activities of data users
Users may have particular data points of interest within a data set, as opposed to the general messaging outlined above. Such low-level user analytic activities are presented in the following table. The taxonomy can also be organized by three poles of activities: retrieving values, finding data points, and arranging data points.
Barriers to effective analysis
Barriers to effective analysis may exist among the analysts performing the data analysis or among the audience. Distinguishing fact from opinion, cognitive biases, and innumeracy are all challenges to sound data analysis.
Confusing fact and opinion
Effective analysis requires obtaining relevant facts to answer questions, support a conclusion or formal opinion, or test hypotheses. Facts by definition are irrefutable, meaning that any person involved in the analysis should be able to agree upon them. For example, in August 2010, the Congressional Budget Office (CBO) estimated that extending the Bush tax cuts of 2001 and 2003 for the 2011–2020 time period would add approximately $3.3 trillion to the national debt. Everyone should be able to agree that indeed this is what CBO reported; they can all examine the report. This makes it a fact. Whether persons agree or disagree with the CBO is their own opinion.
As another example, the auditor of a public company must arrive at a formal opinion on whether financial statements of publicly traded corporations are "fairly stated, in all material respects". This requires extensive analysis of factual data and evidence to support their opinion. When making the leap from facts to opinions, there is always the possibility that the opinion is erroneous.
Cognitive biases
There are a variety of cognitive biases that can adversely affect analysis. For example, confirmation bias is the tendency to search for or interpret information in a way that confirms one's preconceptions. In addition, individuals may discredit information that does not support their views.
Analysts may be trained specifically to be aware of these biases and how to overcome them. In his book Psychology of Intelligence Analysis, retired CIA analyst Richards Heuer wrote that analysts should clearly delineate their assumptions and chains of inference and specify the degree and source of the uncertainty involved in the conclusions. He emphasized procedures to help surface and debate alternative points of view.
Innumeracy
Effective analysts are generally adept with a variety of numerical techniques. However, audiences may not have such literacy with numbers or numeracy; they are said to be innumerate. Persons communicating the data may also be attempting to mislead or misinform, deliberately using bad numerical techniques.
For example, whether a number is rising or falling may not be the key factor. More important may be the number relative to another number, such as the size of government revenue or spending relative to the size of the economy (GDP) or the amount of cost relative to revenue in corporate financial statements. This numerical technique is referred to as normalization or common-sizing. There are many such techniques employed by analysts, whether adjusting for inflation (i.e., comparing real vs. nominal data) or considering population increases, demographics, etc. Analysts apply a variety of techniques to address the various quantitative messages described in the section above.
Analysts may also analyze data under different assumptions or scenario. For example, when analysts perform financial statement analysis, they will often recast the financial statements under different assumptions to help arrive at an estimate of future cash flow, which they then discount to present value based on some interest rate, to determine the valuation of the company or its stock. Similarly, the CBO analyzes the effects of various policy options on the government's revenue, outlays and deficits, creating alternative future scenarios for key measures.
Other topics
Smart buildings
A data analytics approach can be used in order to predict energy consumption in buildings. The different steps of the data analysis process are carried out in order to realise smart buildings, where the building management and control operations including heating, ventilation, air conditioning, lighting and security are realised automatically by miming the needs of the building users and optimising resources like energy and time.
Analytics and business intelligence
Analytics is the "extensive use of data, statistical and quantitative analysis, explanatory and predictive models, and fact-based management to drive decisions and actions." It is a subset of business intelligence, which is a set of technologies and processes that uses data to understand and analyze business performance to drive decision-making .
Education
In education, most educators have access to a data system for the purpose of analyzing student data. These data systems present data to educators in an over-the-counter data format (embedding labels, supplemental documentation, and a help system and making key package/display and content decisions) to improve the accuracy of educators' data analyses.
Practitioner notes
This section contains rather technical explanations that may assist practitioners but are beyond the typical scope of a Wikipedia article.
Initial data analysis
The most important distinction between the initial data analysis phase and the main analysis phase, is that during initial data analysis one refrains from any analysis that is aimed at answering the original research question. The initial data analysis phase is guided by the following four questions:
Quality of data
The quality of the data should be checked as early as possible. Data quality can be assessed in several ways, using different types of analysis: frequency counts, descriptive statistics (mean, standard deviation, median), normality (skewness, kurtosis, frequency histograms), normal imputation is needed.
Analysis of extreme observations: outlying observations in the data are analyzed to see if they seem to disturb the distribution.
Comparison and correction of differences in coding schemes: variables are compared with coding schemes of variables external to the data set, and possibly corrected if coding schemes are not comparable.
Test for common-method variance.
The choice of analyses to assess the data quality during the initial data analysis phase depends on the analyses that will be conducted in the main analysis phase.
Quality of measurements
The quality of the measurement instruments should only be checked during the initial data analysis phase when this is not the focus or research question of the study. One should check whether structure of measurement instruments corresponds to structure reported in the literature.
There are two ways to assess measurement quality:
Confirmatory factor analysis
Analysis of homogeneity (internal consistency), which gives an indication of the reliability of a measurement instrument. During this analysis, one inspects the variances of the items and the scales, the Cronbach's α of the scales, and the change in the Cronbach's alpha when an item would be deleted from a scale
Initial transformations
After assessing the quality of the data and of the measurements, one might decide to impute missing data, or to perform initial transformations of one or more variables, although this can also be done during the main analysis phase.
Possible transformations of variables are:
Square root transformation (if the distribution differs moderately from normal)
Log-transformation (if the distribution differs substantially from normal)
Inverse transformation (if the distribution differs severely from normal)
Make categorical (ordinal / dichotomous) (if the distribution differs severely from normal, and no transformations help)
Did the implementation of the study fulfill the intentions of the research design?
One should check the success of the randomization procedure, for instance by checking whether background and substantive variables are equally distributed within and across groups. If the study did not need or use a randomization procedure, one should check the success of the non-random sampling, for instance by checking whether all subgroups of the population of interest are represented in sample.Other possible data distortions that should be checked are:
dropout (this should be identified during the initial data analysis phase)
Item non-response (whether this is random or not should be assessed during the initial data analysis phase)
Treatment quality (using manipulation checks).
Characteristics of data sample
In any report or article, the structure of the sample must be accurately described. It is especially important to exactly determine the structure of the sample (and specifically the size of the subgroups) when subgroup analyses will be performed during the main analysis phase.The characteristics of the data sample can be assessed by looking at:
Basic statistics of important variables
Scatter plots
Correlations and associations
Cross-tabulations
Final stage of the initial data analysis
During the final stage, the findings of the initial data analysis are documented, and necessary, preferable, and possible corrective actions are taken.Also, the original plan for the main data analyses can and should be specified in more detail or rewritten. In order to do this, several decisions about the main data analyses can and should be made:
In the case of non-normals: should one transform variables; make variables categorical (ordinal/dichotomous); adapt the analysis method?
In the case of missing data: should one neglect or impute the missing data; which imputation technique should be used?
In the case of outliers: should one use robust analysis techniques?
In case items do not fit the scale: should one adapt the measurement instrument by omitting items, or rather ensure comparability with other (uses of the) measurement instrument(s)?
In the case of (too) small subgroups: should one drop the hypothesis about inter-group differences, or use small sample techniques, like exact tests or bootstrapping?
In case the randomization procedure seems to be defective: can and should one calculate propensity scores and include them as covariates in the main analyses?
Analysis
Several analyses can be used during the initial data analysis phase:
Univariate statistics (single variable)
Bivariate associations (correlations)
Graphical techniques (scatter plots)
It is important to take the measurement levels of the variables into account for the analyses, as special statistical techniques are available for each level:
Nominal and ordinal variables
Frequency counts (numbers and percentages)
Associations
circumambulations (crosstabulations)
hierarchical loglinear analysis (restricted to a maximum of 8 variables)
loglinear analysis (to identify relevant/important variables and possible confounders)
Exact tests or bootstrapping (in case subgroups are small)
Computation of new variables
Continuous variables
Distribution
Statistics (M, SD, variance, skewness, kurtosis)
Stem-and-leaf displays
Box plots
Nonlinear analysis
Nonlinear analysis is often necessary when the data is recorded from a nonlinear system. Nonlinear systems can exhibit complex dynamic effects including bifurcations, chaos, harmonics and subharmonics that cannot be analyzed using simple linear methods. Nonlinear data analysis is closely related to nonlinear system identification.
Main data analysis
In the main analysis phase, analyses aimed at answering the research question are performed as well as any other relevant analysis needed to write the first draft of the research report.
Exploratory and confirmatory approaches
In the main analysis phase, either an exploratory or confirmatory approach can be adopted. Usually the approach is decided before data is collected. In an exploratory analysis no clear hypothesis is stated before analysing the data, and the data is searched for models that describe the data well. In a confirmatory analysis clear hypotheses about the data are tested.
Exploratory data analysis should be interpreted carefully. When testing multiple models at once there is a high chance on finding at least one of them to be significant, but this can be due to a type 1 error. It is important to always adjust the significance level when testing multiple models with, for example, a Bonferroni correction. Also, one should not follow up an exploratory analysis with a confirmatory analysis in the same dataset. An exploratory analysis is used to find ideas for a theory, but not to test that theory as well. When a model is found exploratory in a dataset, then following up that analysis with a confirmatory analysis in the same dataset could simply mean that the results of the confirmatory analysis are due to the same type 1 error that resulted in the exploratory model in the first place. The confirmatory analysis therefore will not be more informative than the original exploratory analysis.
Stability of results
It is important to obtain some indication about how generalizable the results are. While this is often difficult to check, one can look at the stability of the results. Are the results reliable and reproducible? There are two main ways of doing that.
Cross-validation. By splitting the data into multiple parts, we can check if an analysis (like a fitted model) based on one part of the data generalizes to another part of the data as well. Cross-validation is generally inappropriate, though, if there are correlations within the data, e.g. with panel data. Hence other methods of validation sometimes need to be used. For more on this topic, see statistical model validation.
Sensitivity analysis. A procedure to study the behavior of a system or model when global parameters are (systematically) varied. One way to do that is via bootstrapping.
Free software for data analysis
Notable free software for data analysis include:
DevInfo – A database system endorsed by the United Nations Development Group for monitoring and analyzing human development.
ELKI – Data mining framework in Java with data mining oriented visualization functions.
KNIME – The Konstanz Information Miner, a user friendly and comprehensive data analytics framework.
Orange – A visual programming tool featuring interactive data visualization and methods for statistical data analysis, data mining, and machine learning.
Pandas – Python library for data analysis.
PAW – FORTRAN/C data analysis framework developed at CERN.
R – A programming language and software environment for statistical computing and graphics.
ROOT – C++ data analysis framework developed at CERN.
SciPy – Python library for scientific computing.
Julia – A programming language well-suited for numerical analysis and computational science.
Reproducible analysis
The typical data analysis workflow involves collecting data, running analyses through various scripts, creating visualizations, and writing reports. However, this workflow presents challenges, including a separation between analysis scripts and data, as well as a gap between analysis and documentation. Often, the correct order of running scripts is only described informally or resides in the data scientist's memory. The potential for losing this information creates issues for reproducibility. To address these challenges, it is essential to have analysis scripts written for automated, reproducible workflows. Additionally, dynamic documentation is crucial, providing reports that are understandable by both machines and humans, ensuring accurate representation of the analysis workflow even as scripts evolve.
International data analysis contests
Different companies or organizations hold data analysis contests to encourage researchers to utilize their data or to solve a particular question using data analysis. A few examples of well-known international data analysis contests are as follows:
Kaggle competition, which is held by Kaggle.
LTPP data analysis contest held by FHWA and ASCE.
| Physical sciences | Science basics | Basics and measurement |
2722105 | https://en.wikipedia.org/wiki/Polarizer | Polarizer | A polarizer or polariser is an optical filter that lets light waves of a specific polarization pass through while blocking light waves of other polarizations. It can filter a beam of light of undefined or mixed polarization into a beam of well-defined polarization, known as polarized light. Polarizers are used in many optical techniques and instruments. Polarizers find applications in photography and LCD technology. In photography, a polarizing filter can be used to filter out reflections.
The common types of polarizers are linear polarizers and circular polarizers. Polarizers can also be made for other types of electromagnetic waves besides visible light, such as radio waves, microwaves, and X-rays.
Linear polarizers
Linear polarizers can be divided into two general categories: absorptive polarizers, where the unwanted polarization states are absorbed by the device, and beam-splitting polarizers, where the unpolarized beam is split into two beams with opposite polarization states. Polarizers which maintain the same axes of polarization with varying angles of incidence are often called Cartesian polarizers, since the polarization vectors can be described with simple Cartesian coordinates (for example, horizontal vs. vertical) independent from the orientation of the polarizer surface. When the two polarization states are relative to the direction of a surface (usually found with Fresnel reflection), they are usually termed s and p. This distinction between Cartesian and s–p polarization can be negligible in many cases, but it becomes significant for achieving high contrast and with wide angular spreads of the incident light.
Absorptive polarizers
Certain crystals, due to the effects described by crystal optics, show dichroism, preferential absorption of light which is polarized in particular directions. They can therefore be used as linear polarizers. The best known crystal of this type is tourmaline. However, this crystal is seldom used as a polarizer, since the dichroic effect is strongly wavelength dependent and the crystal appears coloured. Herapathite is also dichroic, and is not strongly coloured, but is difficult to grow in large crystals.
A Polaroid polarizing filter functions similarly on an atomic scale to the wire-grid polarizer. It was originally made of microscopic herapathite crystals. Its current H-sheet form is made from polyvinyl alcohol (PVA) plastic with an iodine doping. Stretching of the sheet during manufacture causes the PVA chains to align in one particular direction. Valence electrons from the iodine dopant are able to move linearly along the polymer chains, but not transverse to them. So incident light polarized parallel to the chains is absorbed by the sheet; light polarized perpendicularly to the chains is transmitted. The durability and practicality of Polaroid makes it the most common type of polarizer in use, for example for sunglasses, photographic filters, and liquid crystal displays. It is also much cheaper than other types of polarizer.
A modern type of absorptive polarizer is made of elongated silver nano-particles embedded in thin (≤0.5 mm) glass plates. These polarizers are more durable, and can polarize light much better than plastic Polaroid film, achieving polarization ratios as high as 100,000:1 and absorption of correctly polarized light as low as 1.5%. Such glass polarizers perform best for long-wavelength infrared light, and are widely used in fiber-optic communication.
Beam-splitting polarizers
Beam-splitting polarizers split the incident beam into two beams of differing linear polarization. For an ideal polarizing beamsplitter these would be fully polarized, with orthogonal polarizations. For many common beam-splitting polarizers, however, only one of the two output beams is fully polarized. The other contains a mixture of polarization states.
Unlike absorptive polarizers, beam splitting polarizers do not need to absorb and dissipate the energy of the rejected polarization state, and so they are more suitable for use with high intensity beams such as laser light. True polarizing beamsplitters are also useful where the two polarization components are to be analyzed or used simultaneously.
Polarization by Fresnel reflection
When light reflects (by Fresnel reflection) at an angle from an interface between two transparent materials, the reflectivity is different for light polarized in the plane of incidence and light polarized perpendicular to it. Light polarized in the plane is said to be p-polarized, while that polarized perpendicular to it is s-polarized. At a special angle known as Brewster's angle, no p-polarized light is reflected from the surface, thus all reflected light must be s-polarized, with an electric field perpendicular to the plane of incidence.
A simple linear polarizer can be made by tilting a stack of glass plates at Brewster's angle to the beam. Some of the s-polarized light is reflected from each surface of each plate. For a stack of plates, each reflection depletes the incident beam of s-polarized light, leaving a greater fraction of p-polarized light in the transmitted beam at each stage. For visible light in air and typical glass, Brewster's angle is about 57°, and about 16% of the s-polarized light present in the beam is reflected for each air-to-glass or glass-to-air transition. It takes many plates to achieve even mediocre polarization of the transmitted beam with this approach. For a stack of 10 plates (20 reflections), about 3% (= (1 − 0.16)20) of the s-polarized light is transmitted. The reflected beam, while fully polarized, is spread out and may not be very useful.
A more useful polarized beam can be obtained by tilting the pile of plates at a steeper angle to the incident beam. Counterintuitively, using incident angles greater than Brewster's angle yields a higher degree of polarization of the transmitted beam, at the expense of decreased overall transmission. For angles of incidence steeper than 80° the polarization of the transmitted beam can approach 100% with as few as four plates, although the transmitted intensity is very low in this case. Adding more plates and reducing the angle allows a better compromise between transmission and polarization to be achieved.
Because their polarization vectors depend on incidence angle, polarizers based on Fresnel reflection inherently tend to produce s–p polarization rather than Cartesian polarization, which limits their use in some applications.
Birefringent polarizers
Other linear polarizers exploit the birefringent properties of crystals such as quartz and calcite. In these crystals, a beam of unpolarized light incident on their surface is split by refraction into two rays. Snell's law holds for both of these rays, the ordinary or o-ray, and the extraordinary or e-ray, with each ray experiencing a different index of refraction (this is called double refraction). In general the two rays will be in different polarization states, though not in linear polarization states except for certain propagation directions relative to the crystal axis.
A Nicol prism was an early type of birefringent polarizer, that consists of a crystal of calcite which has been split and rejoined with Canada balsam. The crystal is cut such that the o- and e-rays are in orthogonal linear polarization states. Total internal reflection of the o-ray occurs at the balsam interface, since it experiences a larger refractive index in calcite than in the balsam, and the ray is deflected to the side of the crystal. The e-ray, which sees a smaller refractive index in the calcite, is transmitted through the interface without deflection. Nicol prisms produce a very high purity of polarized light, and were extensively used in microscopy, though in modern use they have been mostly replaced with alternatives such as the Glan–Thompson prism, Glan–Foucault prism, and Glan–Taylor prism. These prisms are not true polarizing beamsplitters since only the transmitted beam is fully polarized.
A Wollaston prism is another birefringent polarizer consisting of two triangular calcite prisms with orthogonal crystal axes that are cemented together. At the internal interface, an unpolarized beam splits into two linearly polarized rays which leave the prism at a divergence angle of 15°–45°. The Rochon and Sénarmont prisms are similar, but use different optical axis orientations in the two prisms. The Sénarmont prism is air spaced, unlike the Wollaston and Rochon prisms. These prisms truly split the beam into two fully polarized beams with perpendicular polarizations. The Nomarski prism is a variant of the Wollaston prism, which is widely used in differential interference contrast microscopy.
Thin film polarizers
Thin-film linear polarizers (also known as TFPN) are glass substrates on which a special optical coating is applied. Either Brewster's angle reflections or interference effects in the film cause them to act as beam-splitting polarizers. The substrate for the film can either be a plate, which is inserted into the beam at a particular angle, or a wedge of glass that is cemented to a second wedge to form a cube with the film cutting diagonally across the center (one form of this is the very common MacNeille cube).
Thin-film polarizers generally do not perform as well as Glan-type polarizers, but they are inexpensive and provide two beams that are about equally well polarized. The cube-type polarizers generally perform better than the plate polarizers. The former are easily confused with Glan-type birefringent polarizers.
Wire-grid polarizers
One of the simplest linear polarizers is the wire-grid polarizer (WGP), which consists of many fine parallel metallic wires placed in a plane. WGPs mostly reflect the non-transmitted polarization and can thus be used as polarizing beam splitters. The parasitic absorption is relatively high compared to most of the dielectric polarizers though much lower than in absorptive polarizers.
Electromagnetic waves that have a component of their electric fields aligned parallel to the wires will induce the movement of electrons along the length of the wires. Since the electrons are free to move in this direction, the polarizer behaves in a similar manner to the surface of a metal when reflecting light, and the wave is reflected backwards along the incident beam (minus a small amount of energy lost to Joule heating of the wire).
For waves with electric fields perpendicular to the wires, the electrons cannot move very far across the width of each wire. Therefore, little energy is reflected and the incident wave is able to pass through the grid. In this case the grid behaves like a dielectric material.
Overall, this causes the transmitted wave to be linearly polarized with an electric field completely perpendicular to the wires. The hypothesis that the waves "slip through" the gaps between the wires is incorrect.
For practical purposes, the separation between wires must be less than the wavelength of the incident radiation. In addition, the width of each wire should be small compared to the spacing between wires. Therefore, it is relatively easy to construct wire-grid polarizers for microwaves, far-infrared, and mid-infrared radiation. For far-infrared optics, the polarizer can be even made as free standing mesh, entirely without transmissive optics. In addition, advanced lithographic techniques can also build very tight pitch metallic grids (typ. 50‒100 nm), allowing for the polarization of visible or infrared light to a useful degree. Since the degree of polarization depends little on wavelength and angle of incidence, they are used for broad-band applications such as projection.
Analytical solutions using rigorous coupled-wave analysis for wire grid polarizers have shown that for electric field components perpendicular to the wires, the medium behaves like a dielectric, and for electric field components parallel to the wires, the medium behaves like a metal (reflective).
Malus' law and other properties
Malus' law (), which is named after Étienne-Louis Malus, says that when a perfect polarizer is placed in a polarized beam of light, the irradiance, I, of the light that passes through is given by
where I0 is the initial intensity and θi is the angle between the light's initial polarization direction and the axis of the polarizer.
A beam of unpolarized light can be thought of as containing a uniform mixture of linear polarizations at all possible angles. Since the average value of is 1/2, the transmission coefficient becomes
In practice, some light is lost in the polarizer and the actual transmission will be somewhat lower than this, around 38% for Polaroid-type polarizers but considerably higher (>49.9%) for some birefringent prism types.
If two polarizers are placed one after another (the second polarizer is generally called an analyzer), the mutual angle between their polarizing axes gives the value of θ in Malus's law. If the two axes are orthogonal, the polarizers are crossed and in theory no light is transmitted, though again practically speaking no polarizer is perfect and the transmission is not exactly zero (for example, crossed Polaroid sheets appear slightly blue in colour because their extinction ratio is better in the red). If a transparent object is placed between the crossed polarizers, any polarization effects present in the sample (such as birefringence) will be shown as an increase in transmission. This effect is used in polarimetry to measure the optical activity of a sample.
Real polarizers are also not perfect blockers of the polarization orthogonal to their polarization axis; the ratio of the transmission of the unwanted component to the wanted component is called the extinction ratio, and varies from around 1:500 for Polaroid to about 1:106 for Glan–Taylor prism polarizers.
In X-ray the Malus' law (relativistic form):
where – frequency of the polarized radiation falling on the polarizer, – frequency of the radiation passes through polarizer, – Compton wavelength of electron, – speed of light in vacuum.
Circular polarizers
Circular polarizers (CPL or circular polarizing filters) can be used to create circularly polarized light or alternatively to selectively absorb or pass clockwise and counter-clockwise circularly polarized light.
They are used as polarizing filters in photography to reduce oblique reflections from non-metallic surfaces, and are the lenses of the 3D glasses worn for viewing some stereoscopic movies (notably, the RealD 3D variety), where the polarization of light is used to differentiate which image should be seen by the left and right eye.
Creating circularly polarized light
There are several ways to create circularly polarized light, the cheapest and most common involves placing a quarter-wave plate after a linear polarizer and directing unpolarized light through the linear polarizer. The linearly polarized light leaving the linear polarizer is transformed into circularly polarized light by the quarter wave plate.
The transmission axis of the linear polarizer needs to be half way (45°) between the fast and slow axes of the quarter-wave plate.
In the arrangement above, the transmission axis of the linear polarizer is at a positive 45° angle relative to the right horizontal and is represented with an orange line. The quarter-wave plate has a horizontal slow axis and a vertical fast axis and they are also represented using orange lines. In this instance the unpolarized light entering the linear polarizer is displayed as a single wave whose amplitude and angle of linear polarization are suddenly changing.
When one attempts to pass unpolarized light through the linear polarizer, only light that has its electric field at the positive 45° angle leaves the linear polarizer and enters the quarter-wave plate. In the illustration, the three wavelengths of unpolarized light represented would be transformed into the three wavelengths of linearly polarized light on the other side of the linear polarizer.
In the illustration toward the right is the electric field of the linearly polarized light just before it enters the quarter-wave plate. The red line and associated field vectors represent how the magnitude and direction of the electric field varies along the direction of travel. For this plane electromagnetic wave, each vector represents the magnitude and direction of the electric field for an entire plane that is perpendicular to the direction of travel. (Refer to these two images in the plane wave article to better appreciate this.)
Light and all other electromagnetic waves have a magnetic field which is in phase with, and perpendicular to, the electric field being displayed in these illustrations.
To understand the effect the quarter-wave plate has on the linearly polarized light it is useful to think of the light as being divided into two components which are at right angles (orthogonal) to each other. Towards this end, the blue and green lines are projections of the red line onto the vertical and horizontal planes respectively and represent how the electric field changes in the direction of those two planes. The two components have the same amplitude and are in phase.
Because the quarter-wave plate is made of a birefringent material, when in the wave plate, the light travels at different speeds depending on the direction of its electric field. This means that the horizontal component which is along the slow axis of the wave plate will travel at a slower speed than the component that is directed along the vertical fast axis. Initially the two components are in phase, but as the two components travel through the wave plate the horizontal component of the light drifts farther behind that of the vertical. By adjusting the thickness of the wave plate one can control how much the horizontal component is delayed relative to vertical component before the light leaves the wave plate and they begin again to travel at the same speed. When the light leaves the quarter-wave plate the rightward horizontal component will be exactly one quarter of a wavelength behind the vertical component making the light left-hand circularly polarized when viewed from the receiver.
At the top of the illustration toward the right is the circularly polarized light after it leaves the wave plate. Directly below it, for comparison purposes, is the linearly polarized light that entered the quarter-wave plate. In the upper image, because this is a plane wave, each vector leading from the axis to the helix represents the magnitude and direction of the electric field for an entire plane that is perpendicular to the direction of travel. All the electric field vectors have the same magnitude indicating that the strength of the electric field does not change. The direction of the electric field however steadily rotates.
The blue and green lines are projections of the helix onto the vertical and horizontal planes respectively and represent how the electric field changes in the direction of those two planes. Notice how the rightward horizontal component is now one quarter of a wavelength behind the vertical component. It is this quarter of a wavelength phase shift that results in the rotational nature of the electric field. When the magnitude of one component is at a maximum the magnitude of the other component is always zero. This is the reason that there are helix vectors which exactly correspond to the maxima of the two components.
In the instance just cited, using the handedness convention used in many optics textbooks, the light is considered left-handed/counter-clockwise circularly polarized. Referring to the accompanying animation, it is considered left-handed because if one points one's left thumb against the direction of travel, ones fingers curl in the direction the electric field rotates as the wave passes a given point in space. The helix also forms a left-handed helix in space. Similarly this light is considered counter-clockwise circularly polarized because if a stationary observer faces against the direction of travel, the person will observe its electric field rotate in the counter-clockwise direction as the wave passes a given point in space.
To create right-handed, clockwise circularly polarized light one simply rotates the axis of the quarter-wave plate 90° relative to the linear polarizer. This reverses the fast and slow axes of the wave plate relative to the transmission axis of the linear polarizer reversing which component leads and which component lags.
In trying to appreciate how the quarter-wave plate transforms the linearly polarized light, it is important to realize that the two components discussed are not entities in and of themselves but are merely mental constructs one uses to help appreciate what is happening. In the case of linearly and circularly polarized light, at each point in space, there is always a single electric field with a distinct vector direction, the quarter-wave plate merely has the effect of transforming this single electric field.
Absorbing and passing circularly polarized light
Circular polarizers can also be used to selectively absorb or pass right-handed or left-handed circularly polarized light. It is this feature which is utilized by the 3D glasses in stereoscopic cinemas such as RealD Cinema. A given polarizer which creates one of the two polarizations of light will pass that same polarization of light when that light is sent through it in the other direction. In contrast it will block light of the opposite polarization.
The illustration above is identical to the previous similar one with the exception that the left-handed circularly polarized light is now approaching the polarizer from the opposite direction and linearly polarized light is exiting the polarizer toward the right.
First note that a quarter-wave plate always transforms circularly polarized light into linearly polarized light. It is only the resulting angle of polarization of the linearly polarized light that is determined by the orientation of the fast and slow axes of the quarter-wave plate and the handedness of the circularly polarized light. In the illustration, the left-handed circularly polarized light entering the polarizer is transformed into linearly polarized light which has its direction of polarization along the transmission axis of the linear polarizer and it therefore passes. In contrast right-handed circularly polarized light would have been transformed into linearly polarized light that had its direction of polarization along the absorbing axis of the linear polarizer, which is at right angles to the transmission axis, and it would have therefore been blocked.
To understand this process, refer to the illustration on the right. It is absolutely identical to the earlier illustration even though the circularly polarized light at the top is now considered to be approaching the polarizer from the left. One can observe from the illustration that the leftward horizontal (as observed looking along the direction of travel) component is leading the vertical component and that when the horizontal component is retarded by one quarter of a wavelength it will be transformed into the linearly polarized light illustrated at the bottom and it will pass through the linear polarizer.
There is a relatively straightforward way to appreciate why a polarizer which creates a given handedness of circularly polarized light also passes that same handedness of polarized light. First, given the dual usefulness of this image, begin by imagining the circularly polarized light displayed at the top as still leaving the quarter-wave plate and traveling toward the left. Observe that had the horizontal component of the linearly polarized light been retarded by a quarter of wavelength twice, which would amount to a full half wavelength, the result would have been linearly polarized light that was at a right angle to the light that entered. If such orthogonally polarized light were rotated on the horizontal plane and directed back through the linear polarizer section of the circular polarizer it would clearly pass through given its orientation. Now imagine the circularly polarized light which has already passed through the quarter-wave plate once, turned around and directed back toward the circular polarizer again. Let the circularly polarized light illustrated at the top now represent that light. Such light is going to travel through the quarter-wave plate a second time before reaching the linear polarizer and in the process, its horizontal component is going to be retarded a second time by one quarter of a wavelength. Whether that horizontal component is retarded by one quarter of a wavelength in two distinct steps or retarded a full half wavelength all at once, the orientation of the resulting linearly polarized light will be such that it passes through the linear polarizer.
Had it been right-handed, clockwise circularly polarized light approaching the circular polarizer from the left, its horizontal component would have also been retarded, however the resulting linearly polarized light would have been polarized along the absorbing axis of the linear polarizer and it would not have passed.
To create a circular polarizer that instead passes right-handed polarized light and absorbs left-handed light, one again rotates the wave plate and linear polarizer 90° relative to each another. It is easy to appreciate that by reversing the positions of the transmitting and absorbing axes of the linear polarizer relative to the quarter-wave plate, one changes which handedness of polarized light gets transmitted and which gets absorbed.
Homogeneous circular polarizer
A homogeneous circular polarizer passes one handedness of circular polarization unaltered and blocks the other handedness. This is similar to the way that a linear polarizer would fully pass one angle of linearly polarized light unaltered, but would fully block any linearly polarized light that was orthogonal to it.
A homogeneous circular polarizer can be created by sandwiching a linear polarizer between two quarter-wave plates. Specifically we take the circular polarizer described previously, which transforms circularly polarized light into linear polarized light, and add to it a second quarter-wave plate rotated 90° relative to the first one.
Generally speaking, and not making direct reference to the above illustration, when either of the two polarizations of circularly polarized light enters the first quarter-wave plate, one of a pair of orthogonal components is retarded by one quarter of a wavelength relative to the other. This creates one of two linear polarizations depending on the handedness the circularly polarized light. The linear polarizer sandwiched between the quarter wave plates is oriented so that it will pass one linear polarization and block the other. The second quarter-wave plate then takes the linearly polarized light that passes and retards the orthogonal component that was not retarded by the previous quarter-wave plate. This brings the two components back into their initial phase relationship, reestablishing the selected circular polarization.
Note that it does not matter in which direction one passes the circularly polarized light.
Circular and linear polarizing filters for photography
Linear polarizing filters were the first types to be used in photography and can still be used for non-reflex and older single-lens reflex cameras (SLRs). However, cameras with through-the-lens metering (TTL) and autofocusing systems – that is, all modern SLR and DSLR – rely on optical elements that pass linearly polarized light. If light entering the camera is already linearly polarized, it can upset the exposure or autofocus systems. Circular polarizing filters cut out linearly polarized light and so can be used to darken skies, improve saturation and remove reflections, but the circular polarized light it passes does not impair through-the-lens systems.
| Technology | Optics | null |
2726466 | https://en.wikipedia.org/wiki/SNi | SNi | In chemistry, Si (substitution nucleophilic internal) refers to a specific, regio-selective but not often encountered reaction mechanism for nucleophilic aliphatic substitution. The name was introduced by Cowdrey et al. in 1937 to label nucleophilic reactions which occur with retention of configuration, but later was employed to describe various reactions that proceed with a similar mechanism.
A typical representative organic reaction displaying this mechanism is the chlorination of alcohols with thionyl chloride, or the decomposition of alkyl chloroformates, the main feature is retention of stereochemical configuration. Some examples for this reaction were reported by Edward S. Lewis and Charles E. Boozer in 1952. Mechanistic and kinetic studies were reported few years later by various researchers.
Thionyl chloride first reacts with the alcohol to form an alkyl chloro sulfite, actually forming an intimate ion pair. The second step is the loss of a sulfur dioxide molecule and its replacement by the chloride, which was attached to the sulphite group. The difference between S1 and Si is actually that the ion pair is not completely dissociated, and therefore no real carbocation is formed, which else would lead to a racemisation.
This reaction type is linked to many forms of neighbouring group participation, for instance the reaction of the sulfur or nitrogen lone pair in sulfur mustard or nitrogen mustard to form the cationic intermediate.
This reaction mechanism is supported by the observation that addition of pyridine to the reaction leads to inversion. The reasoning behind this finding is that pyridine reacts with the intermediate sulfite replacing chlorine. The dislodged chlorine has to resort to nucleophilic attack from the rear as in a regular nucleophilic substitution.
In the complete picture for this reaction the sulfite reacts with a chlorine ion in a standard S2 reaction with inversion of configuration. When the solvent is also a nucleophile such as dioxane two successive S2 reactions take place and the stereochemistry is again retention. With standard S1 reaction conditions the reaction outcome is retention via a competing Si mechanism and not racemization and with pyridine added the result is again inversion.
| Physical sciences | Organic reactions | Chemistry |
27394479 | https://en.wikipedia.org/wiki/Gyrfalcon | Gyrfalcon | The gyrfalcon ( or ) (), also abbreviated as gyr, is a bird of prey from the genus Falco (falcons and kestrels) and the largest species of the family Falconidae. A high-latitude species, the gyrfalcon breeds on the Arctic coasts and tundra, the islands of northern North America and the Eurosiberian region, where it is mainly a resident species. Some gyrfalcons disperse more widely after the breeding season or in winter, and individual vagrancy can take birds for long distances. Its plumage varies with location, with birds being coloured from all-white to dark brown. These colour variations are called morphs. Like other falcons, it shows sexual dimorphism, with the female much larger than the male.
For centuries, the gyrfalcon has been valued as a hunting bird. Typical prey includes the ptarmigan and waterfowl, which it may attack in flight; and it also hunts fish and small mammals.
Taxonomy and etymology
The gyrfalcon was formally described by Swedish naturalist Carl Linnaeus in 1758 in the tenth edition of his Systema Naturae under its current binomial name Falco rusticolus. The genus name is the Late Latin term for a falcon, , from a sickle, referencing the talons of the bird. The species name is from the Latin , a countryside-dweller, from , "country" and , "to dwell". The bird's common name comes from French ; in Medieval Latin, it is . The first part of the word may come from Old High German ( modern German ; ultimately from Proto-Germanic ("greed")) for "vulture", referring to its size in comparison with other falcons; or from the Latin for "circle" or "curved path", in turn from the Ancient Greek , , meaning "circle" – from the species' circling as it searches for prey, distinct from the hunting of other falcons in its range. The male gyrfalcon is called a gyrkin in falconry.
Description
The gyrfalcon is the largest falcon in the world, being about the same size as the largest buteos but probably slightly heavier. Males are long, weigh , with average weights reported as and have a wingspan from . Females are bulkier and larger, at long, wingspan, and of weight, with average weights of . An outsized female from eastern Siberia was found to have scaled . Among standard measurements, the wing chord is , the tail is , the culmen is and the tarsus is . The gyrfalcon is larger, broader-winged and longer-tailed than the peregrine falcon, which it is known to compete with (and occasionally hunt). It differs from the buzzard in general structure, having pointed wings.
The gyrfalcon is a very polymorphic species, so its plumage varies greatly. The archetypal morphs are called "white", "silver", "brown", and "black", though they can be coloured on a spectrum from all-white to very dark. The brown form of the gyrfalcon is distinguished from the peregrine by the cream streaking on the nape and crown and by the absence of a well-defined malar stripe and cap. The black morph is similar but has a strongly black-spotted underside, rather than finely barred as in the peregrine and the brown-morph gyrfalcon. White form gyrfalcons are the only predominantly white falcons. Silver gyrfalcons resemble a light grey lanner falcon of larger size. The species shows no sex-based colour differences; juveniles are darker and browner than adults.
The black color seems to be sex-linked and to occur mostly in females; it proved difficult for breeders to get males darker than the dark side of slate grey. A color variety that arose in captive breeding is "black chick".
Systematics and evolution
The gyrfalcon is a member of the hierofalcon complex. In this group, ample evidence indicates hybridisation and incomplete lineage sorting, which confounds analyses of DNA sequence data to a massive extent. The radiation of the entire living diversity of hierofalcons took place around the Eemian Stage at the start of the Late Pleistocene. It represents lineages that expanded into the Holarctic and adapted to local conditions; this is in contrast to less northerly populations of northeastern Africa (where the radiation probably originated) that evolved into the saker falcon. Previous beliefs held that gyrfalcons hybridized with sakers in the Altai Mountains, and this gene flow contributed to the genetic lineage of the Altai falcon. However, recent genetic research has not found distinct genetic clusters differentiating Altai falcons from eastern saker falcons (Falco cherrug milvipes), nor evidence supporting the hybridization theory. Instead, this research suggests that gyrfalcons may have evolved from eastern saker falcons, explaining their close genetic relationship.
Some correlation exists between locality and colour morph. Greenland gyrfalcons are lightest, with white plumage flecked with grey on the back and wings being most common. Other subpopulations have varying amounts of the darker morphs: the Icelandic birds tend towards pale, whereas the Eurasian populations are considerably darker and typically incorporate no white birds. Natural separation into regional subspecies is prevented by gyrfalcons' habit of flying long distances whilst exchanging alleles between subpopulations; thus, the allele distributions for the color polymorphism form clines and in darker birds of unknown origin, theoretically any allele combination might be present. For instance, a mating of a pair of captive gyrfalcons is documented to have produced a clutch of four young: one white, one silver, one brown, and one black. Molecular work suggests plumage color is associated with the melanocortin 1 receptor gene (MC1R), where a nonsynonymous point substitution was perfectly associated with the white/melanic polymorphism.
In general, geographic variation follows Bergmann's rule for size and the demands of crypsis for plumage coloration. Several subspecies have been named according to perceived differences between populations but none of these are consistent and thus no living subspecies are currently accepted. The Icelandic population described as F. r. islandus is perhaps the most distinct. The predominantly white Arctic forms are parapatric and seamlessly grade into the subarctic populations. The Icelandic types are presumed to have less gene flow with their neighbors; they show less variation in plumage colors. Comprehensive phylogeographic studies to determine the proper status of the Icelandic population have yet to be performed.
A population genetic study, however, identified the Iceland population as genetically unique relative to other sampled populations in both eastern and western Greenland, Canada, Alaska, and Norway. Further, within Greenland, differing levels of gene flow between western and eastern sampling locations were identified, with apparent asymmetric dispersal in western Greenland from north to south. This dispersal bias is in agreement with the distribution of plumage colour variants with white gyrfalcons in much higher proportion in north Greenland. Although further work is required to determine the ecological factors contributing to these distributions relative to plumage differences, a study using demographic data suggested that plumage color distribution in Greenland may be influenced by nesting chronology with white individuals and pairs laying eggs earlier in the breeding season and producing more offspring.
Swarth's gyrfalcon
A paleosubspecies, Falco rusticolus swarthi, existed during the Late Pleistocene (125,000 to 13,000 years ago). Fossils found in Little Box Elder Cave (Converse County, Wyoming), Dark Canyon Cave (Eddy County, New Mexico), and McKittrick, California were initially described as Falco swarthi ("Swarth falcon" or more properly "Swarth's gyrfalcon") on account of their distinct size. They have meanwhile proven to be largely inseparable from those of living gyrfalcons, except for being somewhat larger.
Swarth's gyrfalcon was on the upper end of the present gyrfalcon's size range, with some stronger females even surpassing it. It seems to have had some adaptations to the temperate semiarid climate that predominated in its range during the last ice age. Ecologically more similar to current Siberian populations (which are generally composed of smaller birds) or to the prairie falcon, this temperate steppe population must have preyed on landbirds and mammals rather than the sea and landbirds which make up much of the American gyrfalcon's diet today.
Ecology
Dietary biology
The gyrfalcon was originally thought to be a bird of tundra and mountains only; however, in June 2011, it was revealed to spend considerable periods during the winter on sea ice far from land. It feeds only on birds and mammals, the latter of which it takes more regularly than many other Falco species. Like other hierofalcons, it usually hunts in a horizontal pursuit, rather than with the peregrine's speedy stoop from a height. Most prey is killed on the ground, whether they are captured there, or if the victim is a flying bird, forced to the ground. The diet is to some extent opportunistic, but a majority breed and hunt coinciding with ptarmigan and seabird colonies. Avian prey can range in size from redpolls around to geese and capercaillies up to in weight, but ptarmigans (Lagopus mutus) and Willow Grouses (L. lagopus) are often chief prey in the tundra. Seabirds such as auks, gulls and seaducks may predominate in coastal areas, and waders and ducks such as mallards (Anas platyrhynchos) on wetlands. Other avian prey include corvids, smaller passerines, doves, and other birds of preys.
Mammalian prey can be locally important, mainly Arctic ground squirrels (Spermophilus parryii) and Arctic hares (Lepus arcticus), and occasionally Norway lemming (Lemmus lemmus) in peak years. Due to the limit of load that they can carry, gyrfalcons mainly take young hares, but both male and female falcons can take down adult hares up in weight and bring dismembered pieces to their nest. Other mammalian prey can include water voles, muskrats, stoats, minks, Arctic fox pups, and rarely also bats. Prey other than birds and mammals are extremely rare, but brown trout (Salmo trutta) have been recorded as prey.
Threat from climate change
In the early 2000s, it was observed that as possible climate change began to temper the Arctic summers, peregrine falcons were expanding their range north to parts of Greenland, and competing with gyrfalcons. Although it is specially adapted for high-Arctic life, and larger than the peregrine, the gyrfalcon is less aggressive and more conflict-averse, and so is less able to compete with peregrines, which can attack and overwhelm the gyrs. However, it remains on the IUCN's Red List with a Conservation Status of Least Concern.
Breeding
The gyrfalcon almost invariably nests on cliff faces. Breeding pairs do not build their own nests, and often use a bare cliff ledge or the abandoned nest of other birds, particularly golden eagles and common ravens. The clutch can range from 1 to 5 eggs, but is usually 2 to 4. The average size of an egg is ; the average weight is . The incubation period averages 35 days, with the chicks hatching at a weight of around . The nestlings are brooded usually for 10 to 15 days and leave the nest at 7 to 8 weeks. At 3 to 4 months of age, the immature gyrfalcons become independent of their parents, though they may associate with their siblings through the following winter.
The only natural predators of gyrfalcons are golden eagles, and even they rarely engage with these formidable falcons. Gyrfalcons have been recorded as aggressively harassing animals that come near their nests, although common ravens are the only predators known to successfully pick off gyrfalcon eggs and hatchlings. Even brown bears have been reportedly dive-bombed. Humans, whether accidentally (automobile collisions or poisoning of carrion to kill mammalian scavengers) or intentionally (through hunting), are the leading cause of death for gyrfalcons. Gyrfalcons that survive into adulthood can live up to 20 years of age.
As F. rusticolus has such a wide range, it is not considered a threatened species by the IUCN. It is not much affected by habitat destruction, but pollution, for instance by pesticides, depressed its numbers in the mid-20th century, and until 1994 it was considered "Near Threatened". Improving environmental standards in developed countries have allowed the birds to make a comeback.
Interaction with humans
The gyrfalcon has long associated with humans, primarily for hunting and in the art of falconry. It is the official bird of Canada's Northwest Territories. The white falcon in the crest of the Icelandic Republic's coat of arms is a variety of gyrfalcon. The white phase gyrfalcon is the official mascot of the United States Air Force Academy.
In the medieval era, the gyrfalcon was considered a royal bird. The geographer and historian Ibn Sa'id al-Maghribi (d. 1286) described certain northern Atlantic islands west of Ireland where these falcons would be brought from, and how the Egyptian Sultan paid 1,000 dinars for each gyrfalcon (or, if it arrived dead, 500 dinars). Due to its rarity and the difficulties involved in obtaining it, in European falconry the gyrfalcon was reserved for kings and nobles; very rarely was a man of lesser rank seen with a gyrfalcon on his fist.
In the 12th century AD China, swan-hunting with gyrfalcons (海東青 hǎidōngqīng in Chinese) obtained from the Jurchen tribes became fashionable among the Khitan nobility. When demand for gyrfalcons exceeded supply, the Liao Emperor imposed a tax payment-in-kind of gyrfalcons on the Jurchen; under the last Liao emperor, tax collectors were entitled to use force to procure sufficient gyrfalcons. This was one cause of the Jurchen rebellion, whose leader Aguda annihilated the Liao empire in 1125, and established the Jin dynasty in its stead.
Falcons are known to be very susceptible to avian influenza. Therefore, an experiment was done with hybrid gyr-saker falcons, which found that five falcons vaccinated with a commercial H5N2 influenza vaccine survived infection with a highly pathogenic H5N1 strain, whereas five unvaccinated falcons died. Thus, both wild and captive gyrfalcons can be protected from bird flu by vaccination.
| Biology and health sciences | Basics | Animals |
4998804 | https://en.wikipedia.org/wiki/Libellula | Libellula | Libellula is a genus of dragonflies, called chasers (in English) or skimmers (in American), in the family Libellulidae. They are distributed throughout the temperate zone of the Northern Hemisphere. Many have showy wing patterns.
Overview
The taxa Ladona (corporals) and Plathemis (whitetails) have been considered as synonyms of Libellula, subgenera, or separate genera by different authorities. Recent phylogenetic analysis has supported their status as either subgenera or full genera.
Species
List of species.
Extant species
Ladona
Plathemis
Fossils
Libellula brodieri†
Libellula calypso†
Libellula doris †
Libellula eusebioi†
Libellula kieseli†
Libellula martini†
Libellula melobasis†
Libellula pannewitziana†
Libellula perse†
Libellula sieboldiana†
Libellula thetis†
Libellula thoe†
Libellula ukrainensis†
| Biology and health sciences | Odonata | Animals |
5000404 | https://en.wikipedia.org/wiki/Field%20line | Field line | A field line is a graphical visual aid for visualizing vector fields. It consists of an imaginary integral curve which is tangent to the field vector at each point along its length. A diagram showing a representative set of neighboring field lines is a common way of depicting a vector field in scientific and mathematical literature; this is called a field line diagram. They are used to show electric fields, magnetic fields, and gravitational fields among many other types. In fluid mechanics, field lines showing the velocity field of a fluid flow are called streamlines.
Definition and description
A vector field defines a direction and magnitude at each point in space. A field line is an integral curve for that vector field and may be constructed by starting at a point and tracing a line through space that follows the direction of the vector field, by making the field line tangent to the field vector at each point. A field line is usually shown as a directed line segment, with an arrowhead indicating the direction of the vector field. For two-dimensional fields the field lines are plane curves; since a plane drawing of a 3-dimensional set of field lines can be visually confusing most field line diagrams are of this type. Since at each point where it is nonzero and finite the vector field has a unique direction, field lines can never intersect, so there is exactly one field line passing through each point at which the vector field is nonzero and finite. Points where the field is zero or infinite have no field line through them, since direction cannot be defined there, but can be the endpoints of field lines.
Since there are an infinite number of points in any region, an infinite number of field lines can be drawn; but only a limited number can be shown on a field line diagram. Therefore which field lines are shown is a choice made by the person or computer program which draws the diagram, and a single vector field may be depicted by different sets of field lines. A field line diagram is necessarily an incomplete description of a vector field, since it gives no information about the field between the drawn field lines, and the choice of how many and which lines to show determines how much useful information the diagram gives.
An individual field line shows the direction of the vector field but not the magnitude. In order to also depict the magnitude of the field, field line diagrams are often drawn so that each line represents the same quantity of flux. Then the density of field lines (number of field lines per unit perpendicular area) at any location is proportional to the magnitude of the vector field at that point. Areas in which neighboring field lines are converging (getting closer together) indicates that the field is getting stronger in that direction.
In vector fields which have nonzero divergence, field lines begin on points of positive divergence (sources) and end on points of negative divergence (sinks), or extend to infinity. For example, electric field lines begin on positive electric charges and end on negative charges. In fields which are divergenceless (solenoidal), such as magnetic fields, field lines have no endpoints; they are either closed loops or are endless.
In physics, drawings of field lines are mainly useful in cases where the sources and sinks, if any, have a physical meaning, as opposed to e.g. the case of a force field of a radial harmonic. For example, Gauss's law states that an electric field has sources at positive charges, sinks at negative charges, and neither elsewhere, so electric field lines start at positive charges and end at negative charges. A gravitational field has no sources, it has sinks at masses, and it has neither elsewhere, gravitational field lines come from infinity and end at masses. A magnetic field has no sources or sinks (Gauss's law for magnetism), so its field lines have no start or end: they can only form closed loops, extend to infinity in both directions, or continue indefinitely without ever crossing itself. However, as stated above, a special situation may occur around points where the field is zero (that cannot be intersected by field lines, because their direction would not be defined) and the simultaneous begin and end of field lines takes place. This situation happens, for instance, in the middle between two identical positive electric point charges. There, the field vanishes and the lines coming axially from the charges end. At the same time, in the transverse plane passing through the middle point, an infinite number of field lines diverge radially. The concomitant presence of the lines that end and begin preserves the divergence-free character of the field in the point.
Note that for this kind of drawing, where the field-line density is intended to be proportional to the field magnitude, it is important to represent all three dimensions. For example, consider the electric field arising from a single, isolated point charge. The electric field lines in this case are straight lines that emanate from the charge uniformly in all directions in three-dimensional space. This means that their density is proportional to , the correct result consistent with Coulomb's law for this case. However, if the electric field lines for this setup were just drawn on a two-dimensional plane, their two-dimensional density would be proportional to , an incorrect result for this situation.
Construction
Given a vector field and a starting point a field line can be constructed iteratively by finding the field vector at that point . The unit tangent vector at that point is: . By moving a short distance along the field direction a new point on the line can be found
Then the field at that point is found and moving a further distance in that direction the next point of the field line is found. At each point the next point can be found by
By repeating this and connecting the points, the field line can be extended as far as desired. This is only an approximation to the actual field line, since each straight segment isn't actually tangent to the field along its length, just at its starting point. But by using a small enough value for , taking a greater number of shorter steps, the field line can be approximated as closely as desired. The field line can be extended in the opposite direction from by taking each step in the opposite direction by using a negative step .
Examples
If the vector field describes a velocity field, then the field lines follow stream lines in the flow. Perhaps the most familiar example of a vector field described by field lines is the magnetic field, which is often depicted using field lines emanating from a magnet.
Divergence and curl
Field lines can be used to trace familiar quantities from vector calculus:
Divergence may be easily seen through field lines, assuming the lines are drawn such that the density of field lines is proportional to the magnitude of the field (see above). In this case, the divergence may be seen as the beginning and ending of field lines. If the vector field is the resultant of radial inverse-square law fields with respect to one or more sources then this corresponds to the fact that the divergence of such a field is zero outside the sources. In a solenoidal vector field (i.e., a vector field where the divergence is zero everywhere), the field lines neither begin nor end; they either form closed loops, or go off to infinity in both directions. If a vector field has positive divergence in some area, there will be field lines starting from points in that area. If a vector field has negative divergence in some area, there will be field lines ending at points in that area.
The Kelvin–Stokes theorem shows that field lines of a vector field with zero curl (i.e., a conservative vector field, e.g. a gravitational field or an electrostatic field) cannot be closed loops. In other words, curl is always present when a field line forms a closed loop. It may be present in other situations too, such as a helical shape of field lines.
Physical significance
While field lines are a "mere" mathematical construction, in some circumstances they take on physical significance. In fluid mechanics, the velocity field lines (streamlines) in steady flow represent the paths of particles of the fluid. In the context of plasma physics, electrons or ions that happen to be on the same field line interact strongly, while particles on different field lines in general do not interact. This is the same behavior that the particles of iron filings exhibit in a magnetic field.
The iron filings in the photo appear to be aligning themselves with discrete field lines, but the situation is more complex. It is easy to visualize as a two-stage-process: first, the filings are spread evenly over the magnetic field but all aligned in the direction of the field. Then, based on the scale and ferromagnetic properties of the filings they damp the field to either side, creating the apparent spaces between the lines that we see. Of course the two stages described here happen concurrently until an equilibrium is achieved. Because the intrinsic magnetism of the filings modifies the field, the lines shown by the filings are only an approximation of the field lines of the original magnetic field. Magnetic fields are continuous, and do not have discrete lines.
| Mathematics | Other | null |
1963231 | https://en.wikipedia.org/wiki/Caridina | Caridina | Caridina is a genus of freshwater atyid shrimp. They are widely found in tropical or subtropical water in Asia, Oceania and Africa. They are filter-feeders and omnivorous scavengers. They range from 0.9 to 9.8 mm (C. cantonensis) to 1.2–7.4 mm (C. serrata) in carapace length.
Taxonomy and species
There is evidence for hybridization between sympatric taxa, requiring care when interpreting molecular phylogenetic analyses that do not use a large number of specimens.
As of March 2022, the Integrated Taxonomic Information System lists the genus Caridina as having 340 species. These include the following species:
Caridina ablepsia Guo, Jiang & Zhang, 1992
Caridina acuta Liang, Chen & W.-X. Li, 2005
Caridina acutirostris Schenkel, 1902
Caridina africana Kingsley, 1883
Caridina alba J. Li & S. Li, 2010
Caridina alphonsi Bouvier, 1919
Caridina amnicolizambezi Richard & Clark, 2009
Caridina amoyensis Liang & Yan, 1977
Caridina angulata Bouvier, 1905
Caridina angustipes Guo & Liang, 2003
Caridina anislaq Cai, Choy & Ng, 2009
Caridina annandalei Kemp, 1918
Caridina apodosis Cai & N. K. Ng, 1999
Caridina appendiculata Jalihal & Shenoy, 1998
Caridina aruensis Roux, 1911
Caridina bakoensis Ng, 1995
Caridina bamaensis Liang & Yan, 1983
Caridina baojingensis Guo, He & Bai, 1992
Caridina barakoma de Mazancourt et al. 2020
Caridina batuan Cai, Choy & Ng, 2009
Caridina belazoniensis Richard & Clark, 2009
Caridina boehmei Klotz & von Rintelen, 2013
Caridina boholensis Cai, Choy & Ng, 2009
Caridina brachydactyla De Man, 1908
Caridina brevidactyla Roux, 1919
Caridina breviata N. K. Ng & Cai, 2000
Caridina brevicarpalis De Man, 1892
Caridina brevispina Liang & Yan, 1986
Caridina bruneiana Choy, 1992
Caridina buehleri Roux, 1934
Caridina buergersi Karge, von Rintelen & Klotz, 2010
Caridina buhi Cai & Shokita, 2006
Caridina bunyonyiensis Richard & Clark, 2005
Caridina burmensis Cai & Ng, 2000
Caridina butonensis Klotz & von Rintelen, 2013
Caridina caerulea von Rintelen & Cai, 2009
Caridina calmani Bouvier, 1919
Caridina camaro Cai, Choy & Ng, 2009
Caridina cantonensis Yü, 1938
Caridina caobangensis S.-Q. Li & Liang, 2002
Caridina carli Roux, 1931
Caridina cavalerieioides Liu & Liang in Liang, 2004
Caridina caverna Liang, Chen & W.-X. Li, 2005
Caridina cavernicola Liang & Zhou, 1993
Caridina cebuensis Cai & Shokita, 2006
Caridina celebensis De Man, 1892
Caridina celestinoi Blanco, 1939
Caridina chauhani Chopra & Tiwari, 1949
Caridina choiseul de Mazancourt et al. 2020
Caridina chishuiensis Cai & Yuan, 1996
Caridina clavipes Guo & Liang, 2003
Caridina clinata Cai, X. Q. Nguyên & Ng, 1999
Caridina cognata De Man, 1915
Caridina confusa Choy & Marshall, 1997
Caridina congoensis Richard & Clark, 2009
Caridina cornuta Liang & Yan, 1986
Caridina costai de Silva, 1982
Caridina crassipes Liang, 1993
Caridina crurispinata Gurney, 1984
Caridina cucphuongensis Đăng, 1980
Caridina curta Liang & Cai, 2000
Caridina demani Roux, 1911
Caridina demenica Cai & Li, 1997
Caridina dennerli von Rintelen & Cai, 2009
Caridina denticulata
Caridina dentifrons N. K. Ng & Cai, 2000
Caridina devaneyi Choy, 1991
Caridina dianchiensis Liang & Yan, 1985
Caridina disjuncta Cai & Liang, 1999
Caridina disparidentata Liang, Yan & Wang, 1984
Caridina ebuneus Richard & Clark, 2009
Caridina edulis Bouvier, 1904
Caridina elisabethae Karge, von Rintelen & Klotz, 2010
Caridina elliptica Cai & Yuan, 1996
Caridina elongapoda Liang & Yan, 1977
Caridina endehensis De Man, 1892
Caridina ensifera Schenkel, 1902
Caridina evae Richard & Clark, 2009
Caridina excavata Kemp, 1913
Caridina excavatoides Johnson, 1961
Caridina fasciata Hung, Chan & Yu, 1993
Caridina fecunda Roux, 1911
Caridina feixiana Cai & Liang, 1999
Caridina fernandoi Arudpragasam & Costa, 1962
Caridina fijiana Choy, 1983
Caridina flavilineata Đăng, 1975
Caridina formosae Hung, Chan & Yu, 1993
Caridina fossarum Heller, 1862
Caridina fusca Klotz, Wowor & von Rintelen, 2021
Caridina gabonensis Roux, 1927
Caridina ghanensis Richard & Clark, 2009
Caridina glaubrechti von Rintelen & Cai, 2009
Caridina glossopoda Liang, Guo & Gao, 1993
Caridina gordonae Richard & Clark, 2005
Caridina gortio Cai & Anker, 2004
Caridina gracilipes De Man, 1892
Caridina gracilirostris De Man, 1892
Caridina gracillima Lanchester, 1901
Caridina grandirostris Stimpson, 1860
Caridina guangxiensis Liang & Zhou, 1993
Caridina gueryi Marquet, Keith & Kalfatak, 2009
Caridina guiyangensis Liang, 2002
Caridina gurneyi Jalihal, Shenoy & Sankolli, 1984
Caridina hainanensis Liang & Yan, 1983
Caridina hanshanensis Tan, 1990
Caridina harmandi Bouvier, 1906
Caridina hodgarti Kemp, 1913
Caridina holthuisi von Rintelen & Cai, 2009
Caridina hongyanensis Cai & Yuan, 1996
Caridina hova Nobili, 1905
Caridina huananensis Liang, 2004
Caridina hubeiensis Liang & S.-Q. Li, 1993
Caridina hunanensis Liang, Guo & Gao, 1993
Caridina imitatrix Holthuis, 1970
Caridina intermedia de Mazancourt et al. 2020
Caridina jalihali Mariappan & Richard, 2006
Caridina jeani Cai, 2010
Caridina jiangxiensis Liang & Zheng, 1985
Caridina johnsoni Cai, Ng & Choy, 2007
Caridina kaombeflutilis Richard & Clark, 2010
Caridina kempi Jalihal, Shenoy & Sankolli, 1984
Caridina kilimae Hilgendorf, 1898
Caridina kunmingensis Z.-Z. Wang & Liang, 2001
Caridina kunnathurensis Richard & Chandran, 1994
Caridina laevis Heller, 1862
Caridina lamiana Holthuis, 1965
Caridina lanceifrons Yu, 1936
Caridina lanceolata Woltereck, 1937
Caridina lanzana Holthuis, 1980
Caridina laoagensis Blanco, 1939
Caridina laroeha Klotz & von Rintelen, 2013
Caridina leclerci Cai & Ng, 2009
Caridina leucosticta Stimpson, 1860
Caridina leytensis Blanco, 1939
Caridina liangi Jiang, Guo & Zhang, 2002
Caridina liaoi Cai, Choy & Ng, 2009
Caridina lilianae Klotz, Wowor & von Rintelen, 2021
Caridina lima Liang, Guo & Gao, 1993
Caridina linduensis Roux, 1904
Caridina lineorostris Richard & Clark, 2009
Caridina lingkonae Woltereck, 1937
Caridina lipalmaria Richard & Clark, 2010
Caridina liui Liang & Yan, 1986
Caridina lobocensis Cai, Choy & Ng, 2009
Caridina loehae Woltereck, 1937
Caridina longa Liang & Yan, 1985
Caridina longiacuta Guo & Wang, 2005
Caridina longicarpus Roux, 1926
Caridina longidigita Cai & Wowor, 2007
Caridina longifrons Cai & Ng, 2007
Caridina longirostris H. Milne-Edwards, 1837
Caridina lovoensis Roth-Woltereck, 1955
Caridina lufengensis Cai & Duan, 1998
Caridina lumilympha Richard & Clark, 2010
Caridina macrodentata Cai & Shokita, 2006
Caridina macrophora Kemp, 1918
Caridina maculata L. Wang, Liang & F. Li, 2008
Caridina maeana de Mazancourt et al. 2020
Caridina maeklongensis Macharoenboon, K., Manonai, V., & Jeratthitikul, E., 2024
Caridina mahalona Cai, Wowor & Choy, 2009
Caridina malawensis Richard & Clark, 2009
Caridina malayensis Cai, Ng & Choy, 2007
Caridina mariae Klotz & von Rintelen, 2014
Caridina marlenae Klotz, Wowor & von Rintelen, 2021
Caridina masapi Woltereck, 1937
Caridina mathiassi Silas & Jayachandran, 2010
Caridina mauritii Bouvier, 1912
Caridina mayamareenae Klotz, Wowor & von Rintelen, 2021
Caridina mccullochi Roux, 1926
Caridina medifolia Cai & Yuan, 1996
Caridina mengae Liang, 1993
Caridina mengaeoides Guo & Suzuki, 1996
Caridina menghaiensis Cai & Dai, 1999
Caridina meridionalis L. Wang, Liang & F. Li, 2008
Caridina mertoni Roux, 1911
Caridina mesofluminis Richard & Clark, 2009
Caridina mindanao Cai & Shokita, 2006
Caridina minidentata Cai & Anker, 2004
Caridina minnanica Liang, 2002
Caridina modiglianii Nobili, 1900
Caridina moeri Roth-Woltereck, 1984
Caridina mongziensis Liang, Yan & Z.-Z. Wang, 1987
Caridina multidentata Stimpson, 1860
Caridina nana de Mazancourt et al. 2020
Caridina nanaoensis Cai & N. K. Ng, 1999
Caridina natalensis Bouvier, 1925
Caridina natarajani Tiwari & R. S. Pillai, 1968
Caridina neglecta Cai & Ng, 2007
Caridina nguyeni S.-Q. Li & Liang, 2002
Caridina nilotica (Roux, 1833)
Caridina norvestica Holthuis, 1965
Caridina novaecaledoniae Roux, 1926
Caridina nudirostris Choy, 1984
Caridina okiamnis Richard & Clark, 2009
Caridina okinawa Cai & Shokita, 2006
Caridina oligospina Liang, Guo & Tang, 1999
Caridina opaensis Roux, 1904
Caridina palawanensis Cai & Shokita, 2006
Caridina panikkari Jalihal, Shenoy & Sankolli, 1984
Caridina papuana Nobili, 1905
Caridina paracornuta Cai & Yuan, 1996
Caridina pareparensis De Man, 1892
Caridina paratypus de Mazancourt et al. 2020
Caridina parvidentata Roux, 1904
Caridina parvirostris De Man, 1892
Caridina parvocula Gurney, 1984
Caridina parvula von Rintelen & Cai, 2009
Caridina paucidentata Wang & Liang, 2005
Caridina paucidentata L.-Q. Wang & Liang, 2005
Caridina pedicultrata Guo & Choy, 1994
Caridina peninsularis Kemp, 1918
Caridina petiti Roux, 1929
Caridina pingi Yü, 1938
Caridina pingioides Yü, 1938
Caridina piokerai de Mazancourt et al. 2020
Caridina pisuku de Mazancourt et al. 2020
Caridina plicata Liang, 2004
Caridina poarae de Mazancourt et al. 2020
Caridina poso Klotz, Wowor & von Rintelen, 2021
Caridina prashadi Tiwari & R. S. Pillai, 1971
Caridina pristis Roux, 1931
Caridina profundicola von Rintelen & Cai, 2009
Caridina propinqua De Man, 1908
Caridina pseudodenticulata Hung, Chan & Yu, 1993
Caridina pseudonilotica Richard & Clark, 2005
Caridina pseudoserrata Đăng & Ðỗ, 2007
Caridina qingyuanensis Guo & He, 2007
Caridina rajadhari Bouvier, 1918
Caridina rangoona Cai & Ng, 2000
Caridina rapaensis Edmondson, 1935
Caridina richtersi Thallwitz, 1892
Caridina roubaudi Bouvier, 1925
Caridina rouxi De Man, 1915
Caridina rubella Fujino & Shokita, 1975
Caridina rubropunctata Đăng & Ðỗ, 2007
Caridina samar Cai & Anker, 2004
Caridina sarasinorum Schenkel, 1902
Caridina schenkeli von Rintelen & Cai, 2009
Caridina semiblepsia Guo, Choy & Gui, 1996
Caridina serrata Stimpson, 1860
Caridina serratirostris De Man, 1892
Caridina shenoyi Jalihal & Sankolli in Jalihal, Shenoy & Sankolli, 1984
Caridina sikipozo de Mazancourt et al. 2020
Caridina shilinica Liang & Cai, 2000
Caridina similis Bouvier, 1904
Caridina simoni Bouvier, 1904
Caridina sinanensis Xu, Li, Zheng & Guo, 2020
Caridina sodenensis Richard & Clark, 2009
Caridina solearipes Guo & De Grave, 1997
Caridina songtaoensis Liang, 2004
Caridina spathulirostris Richters, 1880
Caridina spelunca Choy, 1996
Caridina sphyrapoda Liang & Zhou, 1993
Caridina spinalifrons Guo & De Grave, 1997
Caridina spinata Woltereck, 1937
Caridina spinipoda Liang, Hong & Yang, 1990
Caridina spinosipes Liang, Guo & Tang, 1999
Caridina spinula Choy & Marshall, 1997
Caridina spongicola Zitzler & Cai, 2006
Caridina steineri Cai, 2005
Caridina striata von Rintelen & Cai, 2009
Caridina subventralis Richard & Clark, 2005
Caridina sulawesi Cai & Ng, 2009
Caridina sumatianica Cai & Yuan, 1996
Caridina sumatrensis De Man, 1892
Caridina sundanella Holthuis, 1978
Caridina susuruflabra Richard & Clark, 2009
Caridina temasek Choy & Ng, 1991
Caridina tenuirostris Woltereck, 1937
Caridina tetrazona Chen, Chen, & Guo, 2020
Caridina thambipillai Johnson, 1961
Caridina thermophila Riek, 1953
Caridina thomasi von Rintelen, Karge & Klotz, 2008
Caridina timorensis De Man, 1893
Caridina togoensis Hilgendorf, 1893
Caridina tonkinensis Bouvier, 1919
Caridina trifasciata Yam & Cai, 2003
Caridina troglodytes Holthuis, 1978
Caridina troglophila Holthuis, 1965
Caridina tumida L. Wang, Liang & F. Li, 2008
Caridina tupaia de Mazancourt, Marquet & Keith, 2019
Caridina turipi de Mazancourt et al. 2020
Caridina typus H. Milne-Edwards, 1837
Caridina uminensis Đăng & Ðỗ, 2007
Caridina umtatensis Richard & Clark, 2009
Caridina unca Gurney, 1984
Caridina valencia Cai, Choy & Ng, 2009
Caridina venusta L. Wang, Liang & F. Li, 2008
Caridina vietriensis Đăng & Ðỗ, 2007
Caridina villadolidi Blanco, 1939
Caridina vitiensis Borradaile, 1899
Caridina weberi De Man, 1892
Caridina williamsi Cai & Ng, 2000
Caridina woltereckae Cai, Wowor & Choy, 2009
Caridina wumingensis Cai & N. K. Ng, 1999
Caridina wyckii (Hickson, 1888)
Caridina xiangnanensis X.-Y. Liu, Guo & Yu, 2006
Caridina xiphias Bouvier, 1925
Caridina yilong Cai & Liang, 1999
Caridina yulinica Cai & N. K. Ng, 1999
Caridina yunnanensis Yü, 1938
Caridina zebra Short, 1993
Caridina zeylanica Arudpragasam & Costa, 1962
Caridina zhejiangensis Liang & Zheng, 1985
Caridina zhongshanica Liang, 2004
A number of phylogenetic studies have questioned the monophyly of Caridina.
Threats and Conservation
As of March 2023, the IUCN Red List lists 56 Caridina species as threatened, with 18 listed as critically endangered, 5 listed as endangered, and 33 listed as vulnerable.
Of these, two (Caridina apodosis and Caridina yilong) are listed as possibly extinct and one (Caridina dennerli) is listed as possibly extinct in the wild.
| Biology and health sciences | Shrimps and prawns | Animals |
1965396 | https://en.wikipedia.org/wiki/Fossa%20%28animal%29 | Fossa (animal) | The fossa (Cryptoprocta ferox; or ; ) is a slender, long-tailed, cat-like mammal that is endemic to Madagascar. It is a member of the carnivoran family Eupleridae.
The fossa is the largest mammalian carnivore on Madagascar and has been compared to a small cougar, as it has convergently evolved many cat-like features. Adults have a head-body length of and weigh between , with the males larger than the females. It has semi-retractable claws (meaning it can extend but not retract its claws fully) and flexible ankles that allow it to climb up and down trees head-first, and also support jumping from tree to tree. A larger relative of the species, Cryptoprocta spelea, probably became extinct before 1400.
The species is widespread, although population densities are usually low. It is found solely in forested habitat, and actively hunts both by day and night. Over 50% of its diet consists of lemurs, the endemic primates found on the island; tenrecs, rodents, lizards, birds, and other animals are also documented as prey. Mating usually occurs in trees on horizontal limbs and can last for several hours. Litters range from one to six pups, which are born altricial (blind and toothless). Infants wean after 4.5 months and are independent after a year. Sexual maturity occurs around three to four years of age, and life expectancy in captivity is 20 years. The fossa is listed as a vulnerable species on the IUCN Red List. It is generally feared by the Malagasy people and is often protected by their fady taboo. The greatest threat to the fossa is habitat destruction.
Its taxonomic classification has been controversial because its physical traits resemble those of cats, yet other traits suggest a close relationship with viverrids. Its classification, along with that of the other Malagasy carnivores, influenced hypotheses about how many times mammalian carnivores have colonized Madagascar. With genetic studies demonstrating that the fossa and all other Malagasy carnivores are most closely related to each other forming a clade, recognized as the family Eupleridae, carnivorans are now thought to have colonized the island once, around 18–20 million years ago.
Etymology
The generic name Cryptoprocta refers to how the animal's anus is hidden by its anal pouch, from the Ancient Greek words crypto- "hidden", and procta "anus". The species name ferox is the Latin adjective "fierce" or "wild".
Its common name comes from the word fosa in Malagasy, an Austronesian language, and some authors have adopted the Malagasy spelling in English.
The word is similar to posa (meaning "cat") in the Iban language (another Austronesian language) from Borneo, and both terms may derive from trade languages from the 1600s. However, an alternative etymology suggests a link to another word that comes from Malay: pusa refers to the Malayan weasel (Mustela nudipes). The Malay word pusa could have become posa for cats in Borneo, while in Madagascar the word could have become fosa to refer to the fossa.
Taxonomy
The fossa was formally described by Edward Turner Bennett on the basis of a specimen from Madagascar sent by Charles Telfair in 1833. The common name is the same as the generic name of the Malagasy civet (Fossa fossana), but they are different species. Because of shared physical traits with viverrids, mongooses, and Felidae, its classification has been controversial. Bennett originally placed the fossa as a type of civet in the family Viverridae, a classification that long remained popular among taxonomists. Its compact braincase, large eye sockets, retractable claws, and specialized carnivorous dentition have also led some taxonomists to associate it with the felids. In 1939, William King Gregory and Milo Hellman placed the fossa in its own subfamily within Felidae, the Cryptoproctinae. George Gaylord Simpson placed it back in Viverridae in 1945, still within its own subfamily, yet conceded it had many cat-like characteristics.
In 1993, Géraldine Veron and François Catzeflis published a DNA hybridization study suggesting that the fossa was more closely related to mongooses (family Herpestidae) than to cats or civets. However, in 1995, Veron's morphological study once again grouped it with Felidae. In 2003, molecular phylogenetic studies using nuclear and mitochondrial genes by Anne Yoder and colleagues showed that all native Malagasy carnivorans share a common ancestry that excludes other carnivores (meaning they form a clade, making them monophyletic) and are most closely related to Asian and African Herpestidae. To reflect these relationships, all Malagasy carnivorans are now placed in a single family, Eupleridae. Within Eupleridae, the fossa is placed in the subfamily Euplerinae along with the falanouc (Eupleres goudoti) and Malagasy civet, but its exact relationships are poorly resolved.
An extinct relative of the fossa was described in 1902 from subfossil remains and recognized as a separate species, Cryptoprocta spelea, in 1935. This species was larger than the living fossa (with a body mass estimate roughly twice as great), but otherwise similar. Across Madagascar, people distinguish two kinds of fossa—a large fosa mainty ("black fossa") and the smaller fosa mena ("reddish fossa")—and a white form has been reported in the southwest. It is unclear whether this is purely folklore or individual variation—related to sex, age or instances of melanism and leucism—or whether there is indeed more than one species of living fossa.
Description
The fossa appears as a diminutive form of a large felid, such as a cougar, but with a slender body and muscular limbs, and a tail nearly as long as the rest of the body. It has a mongoose-like head, relatively longer than that of a cat, although with a muzzle that is broad and short, and with large but rounded ears. It has medium brown eyes set relatively wide apart with pupils that contract to slits. Like many carnivorans that hunt at night, its eyes reflect light; the reflected light is orange in hue.
Its head-body length is and its tail is long. There is some sexual dimorphism, with adult males (weighing ) being larger than females (). Smaller individuals are typically found north and east on Madagascar, larger ones to the south and west. Unusually large individuals weighing up to have been reported, but there is some doubt as to the reliability of the measurements. The fossa can smell, hear, and see well. It is a robust animal and illnesses are rare in captive fossas.
Both males and females have short, straight fur that is relatively dense and without spots or patterns. Both sexes are generally a reddish-brown dorsally and colored a dirty cream ventrally. When in rut, they may have an orange coloration to their abdomen from a reddish substance secreted by a chest gland, but this has not been consistently observed by all researchers. The tail tends to be lighter in coloration than the sides. Juveniles are either gray or nearly white.
Several of the animal's physical features are adaptions to climbing through trees. It uses its tail to assist balance and has semi-retractable claws that it uses to climb trees in its search for prey. It has semiplantigrade feet, switching between a plantigrade-like gait (when arboreal) and a digitigrade-like one (when terrestrial). The soles of its paws are nearly bare and covered with strong pads. The fossa has very flexible ankles that allow it to readily grasp tree trunks so as to climb up or down trees head first or to leap to another tree. Captive juveniles have been known to swing upside down by their hindfeet from knotted ropes.
The fossa has several scent glands, although the glands are less developed in females. Like herpestids it has a perianal skin gland inside an anal sac which surrounds the anus like a pocket. The pocket opens to the exterior with a horizontal slit below the tail. Other glands are located near the penis or vagina, with the penile glands emitting a strong odor. Like the herpestids, it has no prescrotal glands.
External genitalia
One of the more peculiar physical features of this species is its external genitalia. The fossa is unique within its family for the shape of its genitalia, which share traits with those of cats and hyenas. The male fossa has an unusually long penis and baculum (penis bone), reaching to between his forelegs when erect, with an average thickness of . The glans extends about halfway down the shaft and is spiny except at the tip. In comparison, the glans of felids is short and spiny, while that of viverrids is smooth and long. The female fossa exhibits transient masculinization, starting at about 1–2 years of age, developing an enlarged, spiny clitoris that resembles a male's penis. The enlarged clitoris is supported by an os clitoridis, which decreases in size as the animal grows. The females do not have a pseudo-scrotum, but they do secrete an orange substance that colors their underparts, much like the secretions of males.
Hormone levels (testosterone, androstenedione, dihydrotestosterone) do not seem to play a part in this transient masculinization, as those levels are the same in masculinized juveniles and non-masculinized adults. It is speculated that the transient masculinization either reduces sexual harassment of juvenile females by adult males, or reduces aggression from territorial females. While females of other mammal species (such as the spotted hyena) have a pseudo-penis, no other is known to diminish in size as the animal grows.
Comparison with related carnivorans
Overall, the fossa has features in common with three different carnivoran families, leading researchers to place it and other members of Eupleridae alternatively in Herpestidae, Viverridae, and Felidae. Felid features are primarily those associated with eating and digestion, including tooth shape and facial portions of the skull, the tongue, and the digestive tract, typical of its exclusively carnivorous diet. The remainder of the skull most closely resembles skulls of genus Viverra, while the general body structure is most similar to that of various members of Herpestidae. The permanent dentition is (three incisors, one canine, three or four premolars, and one molar on each side of both the upper and lower jaws), with the deciduous formula being similar but lacking the fourth premolar and the molar.
The fossa has a large, prominent rhinarium similar to that of viverrids, but has comparatively larger, round ears, almost as large as those of a similarly sized felid. Its facial vibrissae (whiskers) are long, with the longest being longer than its head. Like some mongoose genera, particularly Galidia (which is now in the fossa's own family, Eupleridae) and Herpestes (of Herpestidae), it has carpal vibrissae as well. Its claws are retractile, but unlike those of Felidae species, they are not hidden in skin sheaths. It has three pairs of nipples (one inguinal, one ventral, and one pectoral).
Habitat and distribution
The fossa has the most widespread geographical range of the Malagasy carnivores, and is generally found in low numbers throughout the island in remaining tracts of forest, preferring pristine undisturbed forest habitat. It is also encountered in some degraded forests, but in lower numbers. Although the fossa is found in all known forest habitats throughout Madagascar, including the western dry deciduous forests, the eastern rainforests, and the southern spiny forests, it is seen more frequently in humid than in dry forests. This may be because the reduced canopy in dry forests provides less shade, and also because the fossa seems to travel more easily in humid forests. It is absent from areas with the heaviest habitat disturbance and, like most of Madagascar's fauna, from the central high plateau of the country.
The fossa has been found across several different elevational gradients in undisturbed portions of protected areas throughout Madagascar. In the Réserve Naturelle Intégrale d'Andringitra, evidence of the fossa has been reported at four different sites ranging from . Its highest known occurrence was reported at ; its presence high on the Andringitra Massif was subsequently confirmed in 1996. Similarly, evidence has been reported of the fossa at the elevational extremes of and in the Andohahela National Park. The presence of the fossa at these locations indicates its ability to adapt to various elevations, consistent with its reported distribution in all Madagascar forest types.
Behavior
The fossa is active during both the day and the night and is considered cathemeral; activity peaks may occur early in the morning, late in the afternoon, and late in the night. The animal generally does not reuse sleeping sites, but females with young do return to the same den. The home ranges of male fossas in Kirindy Forest are up to large, compared to for females. These ranges overlap—by about 30 percent according to data from the eastern forests—but females usually have separated ranges. Home ranges grow during the dry season, perhaps because less food and water is available. In general, radio-collared fossas travel between per day, although in one reported case a fossa was observed moving a straight-line distance of in 16 hours.
The animal's population density appears to be low: in Kirindy Forest, where it is thought to be common, its density has been estimated at one animal per in 1998. Another study in the same forest between 1994 and 1996 using the mark and recapture method indicated a population density of one animal per and one adult per .
Except for mothers with young and occasional observations of pairs of males, animals are usually found alone, so that the species is considered solitary. A 2009 publication, however, reported a detailed observation of cooperative hunting, wherein three male fossas hunted a sifaka (Propithecus verreauxi) for 45 minutes, and subsequently shared the prey. This behavior may be a vestige of cooperative hunting that would have been required to take down larger recently extinct lemurs.
Fossas communicate using sounds, scents, and visual signals. Vocalizations include purring, a threatening call, and a call of fear, consisting of "repeated loud, coarse inhalations and gasps of breath". A long, high yelp may function to attract other fossas. Females mew during mating and males produce a sigh when they have found a female. Throughout the year, animals produce long-lasting scent marks on rocks, trees, and the ground using glands in the anal region and on the chest. They also communicate using face and body expression, but the significance of these signals is uncertain. The animal is aggressive only during mating, and males in particular fight boldly. After a short fight, the loser flees and is followed by the winner for a short distance. In captivity, fossas are usually not aggressive and sometimes even allow themselves to be stroked by a zookeeper, but adult males in particular may try to bite.
Diet
The fossa is a carnivore that hunts small to medium-sized animals. One of eight carnivorous species endemic to Madagascar, the fossa is the island's largest surviving endemic terrestrial mammal and the only predator capable of preying upon adults of all extant lemur species, the largest of which can weigh as much as 90 percent of the weight of the average fossa. Although it is the predominant predator of lemurs, reports of its dietary habits demonstrate a wide variety of prey selectivity and specialization depending on habitat and season; diet does not vary by sex. While the fossa is thought to be a lemur specialist in Ranomafana National Park, its diet is more variable in other rain forest habitats.
The diet of the fossa in the wild has been studied by analyzing their distinctive scats, which resemble gray cylinders with twisted ends and measure long by thick. Scat collected and analyzed from both Andohahela and Andringitra contained lemur matter and rodents. Eastern populations in Andringitra incorporate the widest recorded variety of prey, including both vertebrates and invertebrates. Vertebrates consumed ranged from reptiles to a wide variety of birds, including both understory and ground birds, and mammals, including insectivores, rodents, and lemurs. Invertebrates eaten by the fossa in the high mountain zone of Andringitra include insects and crabs. One study found that vertebrates comprised 94% of the diet of fossas, with lemurs comprising over 50%, followed by tenrecs (9%), lizards (9%), and birds (2%). Seeds, which comprised 5% of the diet, may have been in the stomachs of the lemurs eaten, or may have been consumed with fruit taken for water, as seeds were more common in the stomach in the dry season. The average prey size varies geographically; it is only in the high mountains of Andringitra, in contrast to in humid forests and over in dry deciduous forests. In a study of fossa diet in the dry deciduous forest of western Madagascar, more than 90% of prey items were vertebrates, and more than 50% were lemurs. The primary diet consisted of approximately six lemur species and two or three spiny tenrec species, along with snakes and small mammals. Generally, the fossa preys upon larger lemurs and rodents in preference to smaller ones.
Prey is obtained by hunting either on the ground or in the trees. During the non-breeding season the fossa hunts individually, but during the breeding season hunting parties may be seen, and these may be pairs or later on mothers and young. One member of the group scales a tree and chases the lemurs from tree to tree, forcing them down to the ground where the other is easily able to capture them. The fossa is known to eviscerate its larger lemur prey, a trait that, along with its distinct scat, helps identify its kills. Long-term observations of the fossa's predation patterns on rainforest sifakas suggest that the fossa hunts in a subsection of their range until prey density is decreased, then moves on. The fossa has been reported to prey on domestic animals, such as goats and small calves, and especially chickens. Food taken in captivity includes amphibians, birds, insects, reptiles, and small- to medium-sized mammals.
This wide variety of prey items taken in various rainforest habitats is similar to the varied dietary composition noted occurring in the dry forests of western Madagascar, as well. As the largest endemic predator on Madagascar, this dietary flexibility combined with a flexible activity pattern has allowed it to exploit a wide variety of niches available throughout the island, making it a potential keystone species for the Madagascar ecosystems.
Breeding
Fossas have a polyandrous mating system. Most of the details of reproduction in wild populations are from the western dry deciduous forests; determining whether certain of these details are applicable to eastern populations will require further field research. Mating typically occurs during September and October, although there are reports of its occurring as late as December, and can be highly conspicuous. In captivity in the Northern Hemisphere, fossas instead mate in the northern spring, from March to July. Intromission usually occurs in trees on horizontal limbs about off the ground. Frequently the same tree is used year after year, with remarkable precision as to the date the season commences. Trees are often near a water source, and have limbs strong enough and wide enough to support the mating pair, about wide. Some mating has been reported on the ground as well.
As many as eight males will be at a mating site, staying in close vicinity to the receptive female. The female seems to choose the male she mates with, and the males compete for the attention of the female with a significant amount of vocalization and antagonistic interactions. The female may choose to mate with several of the males, and her choice of mate does not seem to have any correlation to the physical appearance of the males. To stimulate the male to mount her, she gives a series of mewling vocalizations. The male mounts from behind, resting his body on her slightly off-center, a position requiring delicate balance; if the female were to stand, the male would have significant difficulty continuing. He places his paws on her shoulders or grasps her around the waist and often licks her neck. Mating may last for nearly three hours. This unusually lengthy mating is due to the physical nature of the male's erect penis, which has backwards-pointing spines along most of its length.
Fossa mating includes a copulatory tie, which may be enforced by the male's spiny penis. The tie is difficult to break if the mating session is interrupted. Copulation with a single male may be repeated several times, with a total mating time of up to fourteen hours, while the male may remain with the female for up to an hour after the mating. A single female may occupy the tree for up to a week, mating with multiple males over that time. Also, other females may take her place, mating with some of the same males as well as others. This mating strategy, whereby the females monopolize a site and maximize the available number of mates, seems to be unique among carnivores. Recent research suggests that this system helps the fossa overcome factors which would normally impede mate-finding, such as low population density and lack of den use.
The birthing of the litter of one to six (typically two to four) takes place in a concealed location, such as an underground den, a termite mound, a rock crevice, or in the hollow of a large tree (particularly those of the genus Commiphora). Contrary to older research, litters are of mixed sexes. Young are born in December or January, making the gestation period 90 days, with the late mating reports indicating a gestational period of about six to seven weeks. The newborns are blind and toothless and weigh no more than . The fur is thin and has been described as gray-brown or nearly white. After about two weeks the cubs' eyes open, they become more active, and their fur darkens to a pearl gray. The cubs do not take solid food until three months old, and do not leave the den until they are 4.5 months old; they are weaned shortly after that. After the first year, the juveniles are independent of their mother. Permanent teeth appear at 18 to 20 months. Physical maturity is reached by about two years of age, but sexual maturity is not attained for another year or two, and the young may stay with their mother until they are fully mature. Lifespan in captivity is up to or past 20 years of age, possibly due to the slow juvenile development.
Human interactions
The fossa has been assessed as "Vulnerable" by the IUCN Red List since 2008, as its population size has probably declined by at least 30 percent between 1987 and 2008; previous assessments have included "Endangered" (2000) and "Insufficiently Known" (1988, 1990, 1994). The species is dependent on forest and thus threatened by the widespread destruction of Madagascar's native forest but is also able to persist in disturbed areas. A suite of microsatellite markers (short segments of DNA that have a repeated sequence) have been developed to help aid in studies of genetic health and population dynamics of both captive and wild fossas. Several pathogens have been isolated from the fossa, some of which, such as anthrax and canine distemper, are thought to have been transmitted by feral dogs or cats. Toxoplasma gondii was reported in a captive fossa in 2013.
Although the species is widely distributed, it is locally rare in all regions, making fossas particularly vulnerable to extinction. The effects of habitat fragmentation increase the risk. For its size, the fossa has a lower than predicted population density, which is further threatened by Madagascar's rapidly disappearing forests and dwindling populations of lemurs, which make up a high proportion of its diet. The loss of the fossa, either locally or completely, could significantly impact ecosystem dynamics, possibly leading to over-grazing by some of its prey species. The total population of the fossa living within protected areas is estimated at less than 2,500 adults, but this may be an overestimate.
Only two protected areas are thought to contain 500 or more adult fossas: Masoala National Park and Midongy-Sud National Park, although these are also thought to be overestimated. Too little population information has been collected for a formal population viability analysis, but estimates suggest that none of the protected areas support a viable population. If this is correct, the extinction of the fossa may take as much as 100 years to occur as the species gradually declines. In order for the species to survive, it is estimated that at least is needed to maintain smaller, short-term viable populations, and at least for populations of 500 adults.
Taboo, known in Madagascar as fady, offers protection for the fossa and other carnivores. In the Marolambo District (part of the Atsinanana region in Toamasina Province), the fossa has traditionally been hated and feared as a dangerous animal. It has been described as "greedy and aggressive", known for taking fowl and piglets, and believed to "take little children who walk alone into the forest". Some do not eat it for fear that it will transfer its undesirable qualities to anyone who consumes it. However, the animal is also taken for bushmeat; a study published in 2009 reported that 57 percent of villages (8 of 14 sampled) in the Makira forest consume fossa meat. The animals were typically hunted using slingshots, with dogs, or most commonly, by placing snare traps on animal paths. Near Ranomafana National Park, the fossa, along with several of its smaller cousins and the introduced small Indian civet (Viverricula indica), are known to "scavenge on the bodies of ancestors", which are buried in shallow graves in the forest. For this reason, eating these animals is strictly prohibited by fady. However, if they wander into villages in search of domestic fowl, they may be killed or trapped. Small carnivore traps have been observed near chicken runs in the village of Vohiparara. Retaliatory killing of fossas in response to their predation on coops is common as well; people with a educational attainment were more likely to dislike fossas, and those that disliked fossas were more likely to report having killed one.
Fossas are occasionally held in captivity in zoos. They first bred in captivity in 1974 in the zoo of Montpellier, France. The next year, at a time when there were only eight fossas in the world's zoos, the Duisburg Zoo in Germany acquired one; this zoo later started a successful breeding program, and most zoo fossas now descend from the Duisburg population. Research on the Duisburg fossas has provided much data about their biology.
The fossa was depicted as an antagonist in the 2005 DreamWorks animated film Madagascar, being referred to as the "foosa", and accurately shown as the lemurs' most feared predator.
| Biology and health sciences | Other carnivora | Animals |
1966070 | https://en.wikipedia.org/wiki/Lead%28II%29%20acetate | Lead(II) acetate | Lead(II) acetate is a white crystalline chemical compound with a slightly sweet taste. Its chemical formula is usually expressed as or , where Ac represents the acetyl group. Like many other lead compounds, it causes lead poisoning. Lead acetate is soluble in water and glycerin. With water it forms the trihydrate, , a colourless or white efflorescent monoclinic crystalline substance.
The substance is used as a reagent to make other lead compounds and as a fixative for some dyes. In low concentrations, it formerly served as the principal active ingredient in progressive types of hair colouring dyes. Lead(II) acetate is also used as a mordant in textile printing and dyeing, and as a drier in paints and varnishes. It was historically used as a sweetener and preservative in wines and in other foods and for cosmetics.
Production
Lead(II) acetate can be made by boiling elemental lead in acetic acid and hydrogen peroxide. This method will also work with lead(II) carbonate or lead(II) oxide.
Lead(II) acetate can also be made by dissolving lead(II) oxide in acetic acid:
Lead(II) acetate can also be made via a single-displacement reaction between copper acetate and lead metal:
Structure
The crystal structure of anhydrous lead(II) acetate has been described as a 2D coordination polymer. In comparison, lead(II) acetate trihydrate's structure is a 1D coordination polymer. In the trihydrate, the Pb2+ ion's coordination sphere consists of nine oxygen atoms belonging to three water molecules, two bidentate acetate groups and two bridging acetate groups. The coordination geometry at Pb is a monocapped square antiprism. The trihydrate thermally decomposes to a hemihydrate, Pb(OAc)2·H2O, and to basic acetates such as Pb4O(OAc)6 and Pb2O(OAc)2.
Uses
Lead acetate is used as a precursor to other lead compounds such as the various carbonate.
Niche and laboratory uses
Lead(II) acetate paper is used to detect the poisonous gas hydrogen sulfide. The gas reacts with lead(II) acetate on the moistened test paper to form a grey precipitate of lead(II) sulfide.
An aqueous solution of lead(II) acetate is a byproduct of the process used in the cleaning and maintenance of stainless steel firearm suppressors (silencers) and compensators when using a 1:1 ratio of hydrogen peroxide and white vinegar (acetic acid). The solution is agitated by the bubbling action of the hydrogen peroxide, with the main reaction being the oxidation of lead by hydrogen peroxide and subsequent dissolution of lead oxide by the acetic acid, which forms lead acetate. Because of its high toxicity, this chemical solution must be appropriately disposed by a chemical processing facility or hazardous materials centre. Alternatively, the solution may be reacted with sulfuric acid to precipitate nearly insoluble lead(II) sulfate. The solid may then be removed by mechanical filtration and is safer to dispose of than aqueous lead acetate.
Historical uses
Sweetener
Like other lead(II) salts, lead(II) acetate has a sweet taste, which led to its historical use as a sugar substitute in both wines and foods.
The ancient Romans, who had few sweeteners besides honey, would boil must (unfiltered grape juice) in lead pots to produce a reduced sugar syrup called defrutum, concentrated again into sapa. This syrup was used to sweeten wine and to sweeten and preserve fruit. It is possible that lead(II) acetate or other lead compounds leaching into the syrup might have caused lead poisoning in those who consumed it. Lead acetate is no longer used in the production of sweeteners because of its recognized toxicity. Legislation prohibiting its use as a wine sweetener was ineffective until decades later, when chemical methods of detecting its presence had been developed.
The earliest confirmed poisoning by lead acetate was that of Pope Clement II, who died in October 1047. A toxicological examination of his remains conducted in the mid-20th century confirmed centuries-old rumors that he had been poisoned with lead sugar. It is not clear whether he was assassinated.
In 1787 painter and biographer Albert Christoph Dies swallowed, by accident, approximately of lead acetate. His recovery from this poison was slow and incomplete. He lived with illnesses until his death in 1822.
Although the use of lead(II) acetate as a sweetener was already illegal at that time, composer Ludwig van Beethoven may have died of lead poisoning caused by wines adulterated with lead acetate (see also Beethoven's liver).
In 1887, 38 hunting horses belonging to Captain William Hollwey Steeds were poisoned in their stables at Clonsilla House, Dublin, Ireland. At least ten of the hunters died. Captain Steeds, an "extensive commission agent", had previously supplied the horses for the Bray and Greystones Coach. It transpired that they had been fed a bran mash that had been sweetened with a toxic lead acetate.
Cosmetics
Lead(II) acetate and white lead have been used in cosmetics throughout history.
It was once used for men's hair colouring products like Grecian Formula. The manufacturer did not remove lead acetate from its product until 2018. Lead acetate has been replaced by bismuth citrate as the progressive colorant. Its use in cosmetics has been banned in Canada by Health Canada since 2005 (effective at the end of 2006) based on tests showing possible carcinogenicity and reproductive toxicity, and it is also banned in the European Union.
Medical uses
Lead(II) acetate solution was a commonly used folk remedy for sore nipples. In modern medicine, for a time, it was used as an astringent, in the form of Goulard's extract, and it has also been used to treat poison ivy.
In the 1850s, Mary Seacole applied lead(II) acetate, among other remedies, against an epidemic of cholera in Panama.
Other historic uses
It was also used in making of slow matches during the Middle Ages. It was made by mixing a natural form of lead(II) oxide called litharge and vinegar.
Sugar of lead was a recommended agent added to linseed oil during heating to produce "boiled" linseed oil, the lead and heat acting to cause the oil to cure faster than raw linseed oil.
Lead(II) acetate ("salt of Saturn") was used to synthesise acetone which was then known as "spirit of Saturn" for being made with the salt of Saturn and thought to be a lead compound in the 17th century.
| Physical sciences | Acetates | Chemistry |
1966840 | https://en.wikipedia.org/wiki/Neisseria%20meningitidis | Neisseria meningitidis | Neisseria meningitidis, often referred to as the meningococcus, is a Gram-negative bacterium that can cause meningitis and other forms of meningococcal disease such as meningococcemia, a life-threatening sepsis. The bacterium is referred to as a coccus because it is round, and more specifically a diplococcus because of its tendency to form pairs.
About 10% of adults are carriers of the bacteria in their nasopharynx. As an exclusively human pathogen, it causes developmental impairment and death in about 10% of cases. It causes the only form of bacterial meningitis known to occur epidemically, mainly in Africa and Asia. It occurs worldwide in both epidemic and endemic form.
N. meningitidis is spread through saliva and respiratory secretions during coughing, sneezing, kissing, chewing on toys and through sharing a source of fresh water. It has also been reported to be transmitted through oral sex and cause urethritis in men. It infects its host cells by sticking to them with long thin extensions called pili and the surface-exposed proteins Opa and Opc and has several virulence factors.
Signs and symptoms
Meningococcus can cause meningitis and other forms of meningococcal disease. It initially produces general symptoms like fatigue, fever, and headache and can rapidly progress to neck stiffness, coma and death in 10% of cases. Petechiae occur in about 50% of cases. Chance of survival is highly correlated with blood cortisol levels, with lower levels prior to steroid administration corresponding with increased patient mortality. Symptoms of meningococcal meningitis are easily confused with those caused by other bacteria, such as Haemophilus influenzae and Streptococcus pneumoniae. Suspicion of meningitis is a medical emergency and immediate medical assessment is recommended. Current guidance in the United Kingdom is that if a case of meningococcal meningitis or septicaemia (infection of the blood) is suspected, intravenous antibiotics should be given and the ill person admitted to the hospital. This means that laboratory tests may be less likely to confirm the presence of Neisseria meningitidis as the antibiotics will dramatically lower the number of bacteria in the body. The UK guidance is based on the idea that the reduced ability to identify the bacteria is outweighed by reduced chance of death.
Septicaemia caused by Neisseria meningitidis has received much less public attention than meningococcal meningitis even though septicaemia has been linked to infant deaths. Meningococcal septicaemia typically causes a purpuric rash, that does not lose its color when pressed with a glass slide ("non-blanching") and does not cause the classical symptoms of meningitis. This means the condition may be ignored by those not aware of the significance of the rash. Septicaemia carries an approximate 50% mortality rate over a few hours from initial onset.
Other severe complications include Waterhouse–Friderichsen syndrome, a massive, usually bilateral, hemorrhage into the adrenal glands caused by fulminant meningococcemia, adrenal insufficiency, and disseminated intravascular coagulation.
Not all instances of a purpura-like rash are due to meningococcal septicaemia; other possible causes, such as idiopathic thrombocytopenic purpura (ITP; a platelet disorder) and Henoch–Schönlein purpura, also need prompt investigation.
Microbiology
N. meningitidis is a Gram-negative diplococcus since it has an outer and inner membranes with a thin layer of peptidoglycan in between. It is 0.6–1.0 micrometers in size. It tests positive for the enzyme cytochrome c oxidase.
Habitat
N. meningitidis is a part of the normal nonpathogenic flora in the nasopharynx of up to 8–25% of adults. It colonizes and infects only humans, and has never been isolated from other animals. This is thought to result from the bacterium's inability to get iron from sources other than human transferrin and lactoferrin.
Subtypes
Disease-causing strains are classified according to the antigenic structure of their polysaccharide capsule. Serotype distribution varies markedly around the world. Among the 13 identified capsular types of N. meningitidis, six (A, B, C, W135, X, and Y) account for most disease cases worldwide. Type A has been the most prevalent in Africa and Asia, but is rare/practically absent in North America. In the United States, serogroup B is the predominant cause of disease and mortality, followed by serogroup C. The multiple subtypes have hindered development of a universal vaccine for meningococcal disease.
Pathogenesis
Virulence
Lipooligosaccharide (LOS) is a component of the outer membrane of N. meningitidis. This acts as an endotoxin and is responsible for septic shock and hemorrhage due to the destruction of red blood cells. Other virulence factors include a polysaccharide capsule which prevents host phagocytosis and aids in evasion of the host immune response. Adhesion is another key virulence strategy to successfully invade host cell. There are several known proteins that are involved in adhesion and invasion, or mediate interactions with specific host cell receptors. These include the Type IV pilin adhesin which mediates attachment of the bacterium to the epithelial cells of the nasopharynx, surface-exposed Opa and Opc proteins which mediate interactions with specific host cell receptors, and NadA which is involved in adhesion.
Pathogenic meningococci that have invaded into the bloodstream must be able to survive in the new niche, this is facilitated by acquisition and utilisation of iron (FetA and Hmbr), resisting intracellular oxidative killing by producing catalase and superoxide dismutase and ability to avoid complement mediated killing (fHbp). Meningococci produce an IgA protease, an enzyme that cleaves IgA class antibodies and thus allows the bacteria to evade a subclass of the humoral immune system.
A hypervirulent strain was discovered in China. Its impact is yet to be determined.
Complement inhibition
Factor H binding protein (fHbp) that is exhibited in N. meningitidis and some commensal species is the main inhibitor of the alternative complement pathway. fHbp protects meningococci from complement-mediated death in human serum experiments, but has also been shown to protect meningococci from antimicrobial peptides in vitro. Factor H binding protein is key to the pathogenesis of N. meningitidis, and is, therefore, important as a potential vaccine candidate. Porins are also an important factor for complement inhibition for both pathogenic and commensal species. Porins are important for nutrient acquisition. Porins are also recognized by TLR2, they bind complement factors (C3b, C4b, factor H, and C4bp (complement factor 4b-binding protein)). Cooperation with pili for CR3-mediated internalization is another function of porins. Ability to translocate into host cells and modulate reactive oxygen species production and apoptosis is made possible by porins, as well. Strains of the same species can express different porins.
Genome
At least 8 complete genomes of Neisseria meningitidis strains have been determined which encode about 2,100 to 2,500 proteins.
The genome of strain MC58 (serogroup B) has 2,272,351 base pairs. When sequenced in 2000, it was found to contain 2158 open reading frames (ORFs). Of these, a biological function was predicted for 1158 (53.7%). There were three major islands of horizontal DNA transfer found. Two encode proteins involved in pathogenicity. The third island only codes for hypothetical proteins. They also found more genes that undergo phase variation than any pathogen then known. Phase variation is a mechanism that helps the pathogen to evade the immune system of the host.
The genome size of strain H44/76 is 2.18 Mb, and encodes 2,480 open reading frames (ORFs), compared to 2.27 Mb and 2,465 ORFs for MC58. Both strains have a GC content of 51.5%. A comparison with MC58 showed that four genes are uniquely present in H44/76 and nine genes are only present in MC58. Of all ORFs in H44/76, 2,317 (93%) show more than 99% sequence identity.
The complete genome sequence of strain NMA510612 (serogroup A) consists of one circular chromosome with a size of 2,188,020 bp, and the average GC content is 51.5%. The chromosome is predicted to possess 4 rRNA operons, 163 insertion elements (IS), 59 tRNAs, and 2,462 ORFs.
There is a public database available for N. meningitidis core genome Multilocus sequence typing (cgMLST). Available at: Neisseria typing
Genetic transformation
Genetic transformation is the process by which a recipient bacterial cell takes up DNA from a neighboring cell and integrates this DNA into the recipient's genome by recombination. In N. meningitidis, DNA transformation requires the presence of short DNA sequences (9–10 mers residing in coding regions) of the donor DNA. These sequences are called DNA uptake sequences (DUSs). Specific recognition of these sequences is mediated by a type IV pilin. In N. meningitidis DUSs occur at a significantly higher density in genes involved in DNA repair and recombination (as well as in restriction-modification and replication) than in other annotated gene groups. The over-representation of DUS in DNA repair and recombination genes may reflect the benefit of maintaining the integrity of the DNA repair and recombination machinery by preferentially taking up genome maintenance genes, that could replace their damaged counterparts in the recipient cell.
N. meningititis colonizes the nasopharyngeal mucosa, which is rich in macrophages. Upon their activation, macrophages produce superoxide (O2−) and hydrogen peroxide (H2O2). Thus N. meningitidis is likely to encounter oxidative stress during its life cycle. Consequently, an important benefit of genetic transformation to N. meningitidis may be the maintenance of the recombination and repair machinery of the cell that removes oxidative DNA damages such as those caused by reactive oxygen. This is consistent with the more general idea that transformation benefits bacterial pathogens by facilitating repair of DNA damages produced by the oxidative defenses of the host during infection.
Meningococci population is extensively diverse genetically, this is due to horizontal gene transfers while in the nasophanrynx. Gene transfer can occur within and between genomes of Neisseria species, and it is the main mechanism of acquiring new traits. This is facilitated by the natural competence of the meningococci to take up foreign DNA. The commensal species of Neisseria can act as a reservoir of genes that can be acquired; for example, this is how capsule switching can occur as a means of hiding from the immune system. An invasive N. meningitidis strain of serogroup C broke out in Nigeria in 2013 – the strain was a new sequence type, ST-10217 determined by multilocus sequence typing. It was determined that a commensal strain of N. meningitidis acquired an 8-kb prophage, the meningococcal disease-associated island (MDAΦ), previously associated with hyper-invasiveness; and the full serogroup C capsule operon, thus becoming a hypervirulent strain. This illustrates how hypervirulent strains can arise from non-pathgenic strains due to the high propensity of gene transfers and DNA uptake by N. meningitidis.
Diagnosis
A small amount of cerebrospinal fluid (CSF) is sent to the laboratory as soon as possible for analysis. The diagnosis is suspected, when Gram-negative diplococci are seen on Gram stain of a centrifuged sample of CSF; sometimes they are located inside white blood cells. The microscopic identification takes around 1–2 hours after specimen arrival in the laboratory.
The gold standard of diagnosis is microbiological isolation of N. meningitidis by growth from a sterile body fluid, which could be CSF or blood. Diagnosis is confirmed when the organism has grown, most often on a chocolate agar plate, but also on Thayer–Martin agar. To differentiate any bacterial growth from other species a small amount of a bacterial colony is Gram stained and tested for oxidase and catalase. Gram negative diplococci that are oxidase and catalase positive are then tested for fermentation of the following carbohydrates: maltose, sucrose, and glucose. N. meningitidis will ferment glucose and maltose. Finally, serology determines the subgroup of the N. meningitidis, which is important for epidemiological surveillance purposes; this may often only be done in specialized laboratories.
The above tests take a minimum of 48–72 hours turnaround time for growing the organism, and up to a week more for serotyping. Growth can and often does fail, either because antibiotics have been given preemptively, or because specimens have been inappropriately transported, as the organism is extremely susceptible to antibiotics and fastidious in its temperature, and growth medium requirements.
Polymerase chain reaction (PCR) tests where available, mostly in industrialized countries, have been increasingly used; PCR can rapidly identify the organism, and works even after antibiotics have been given.
Prevention
All recent contacts of the infected patient over the seven days before onset should receive medication to prevent them from contracting the infection. This especially includes young children and their child caregivers or nursery-school contacts, as well as anyone who had direct exposure to the patient through kissing, sharing utensils, or medical interventions such as mouth-to-mouth resuscitation. Anyone who frequently ate, slept or stayed at the patient's home during the seven days before the onset of symptom, or those who sat beside the patient on an airplane flight or classroom for eight hours or longer, should also receive chemoprophylaxis. The agent of choice is usually oral rifampicin for a few days.
Receiving a dose of the meningococcal vaccine before traveling to a country in the "meningitis belt" or having a booster meningitis vaccine, normally five years apart could prevent someone from getting an infection from the pathogen.
Vaccination
United States
A number of vaccines are available in the U.S. to prevent meningococcal disease. Some of the vaccines cover serogroup B, while others cover A, C, W, and Y. The Centers for Disease Control and Prevention (CDC) recommends all teenagers receive MenACWY vaccine and booster, with optional MenB. MenACWY and MenB are also recommended for people of other ages with various medical conditions and social risk factors.
A meningococcal polysaccharide vaccine (MPSV4) has been available since the 1970s and is the only meningococcal vaccine licensed for people older than 55. MPSV4 may be used in people 2–55 years old if the MCV4 vaccines are not available or contraindicated. Two meningococcal conjugate vaccines (MCV4) are licensed for use in the U.S. The first conjugate vaccine was licensed in 2005, the second in 2010. Conjugate vaccines are the preferred vaccine for people 2 through 55 years of age. It is indicated in those with impaired immunity, such as nephrotic syndrome or splenectomy.
In June 2012, the U.S. Food and Drug Administration (FDA) approved a combination vaccine against two types of meningococcal diseases and Hib disease for infants and children 6 weeks to 18 months old. The vaccine, Menhibrix, was designed to prevent disease caused by Neisseria meningitidis serogroups C and Y, and Haemophilus influenzae type b (Hib). It was the first meningococcal vaccine that could be given to infants as young as six weeks old.
In October 2014 the FDA approved the first vaccine effective against serogroup B, named Trumenba, for use in 10- to 25-year-old individuals.
Africa
In 2010, the Meningitis Vaccine Project introduced a vaccine called MenAfriVac in the African meningitis belt. It was made by generic drug maker Serum Institute of India and cost 50 U.S. cents per injection. Beginning in Burkina Faso in 2010, it has been given to 215 million people across Benin, Cameroon, Chad, Ivory Coast, Ethiopia, Ghana, Mali, Niger, Mauritania, Nigeria, Senegal, Sudan, Togo and Gambia. The vaccination campaign has resulted in near-elimination of serogroup A meningitis from the participating countries.
Treatment
Persons with confirmed N. meningitidis infection should be hospitalized immediately for treatment with antibiotics. Because meningococcal disease can disseminate very rapidly, a single dose of intramuscular antibiotic is often given at the earliest possible opportunity, even before hospitalization, if disease symptoms look suspicious enough. Third-generation cephalosporin antibiotics (i.e. cefotaxime, ceftriaxone) should be used to treat a suspected or culture-proven meningococcal infection before antibiotic susceptibility results are available. Clinical practice guidelines endorse empirical treatment in the event a lumbar puncture to collect cerebrospinal fluid (CSF) for laboratory testing cannot first be performed. Antibiotic treatment may affect the results of microbiology tests, but a diagnosis may be made on the basis of blood-cultures and clinical examination.
Epidemiology
N. meningitidis is a major cause of illness, developmental impairment and death during childhood in industrialized countries and has been responsible for epidemics in Africa and in Asia. Every year, about 2,500 to 3,500 people become infected with N. meningitidis in the US, with a frequency of about 1 in 100,000. Children younger than five years are at greatest risk, followed by teenagers of high school age. Rates in the African meningitis belt were as high as 1 in 1,000 to 1 in 100 before introduction of a vaccine in 2010. The incidence of meningococcal disease is highest among infants (children younger than one year old) whose immune system is relatively immature. In industrialized countries there is a second peak of incidence in young adults, who are congregating closely, living in dormitories or smoking. Vaccine development is ongoing.
It is spread through saliva and other respiratory secretions during coughing, sneezing, kissing, and chewing on toys. Inhalation of respiratory droplets from a carrier which may be someone who is themselves in the early stages of disease can transmit the bacteria. Close contact with a carrier is the predominant risk factor. Other risk factors include a weakened general or local immune response, such as a recent upper respiratory infection, smoking, and complement deficiency. The incubation period is short, from 2 to 10 days. In susceptible individuals, N. meningitidis may invade the bloodstream and cause a systemic infection, sepsis, disseminated intravascular coagulation, breakdown of circulation, and septic shock.
History
In 1884 Ettore Marchiafava and Angelo Celli first observed the bacterium inside cells in the cerebral spinal fluid (CSF). In 1887 Anton Weichselbaum isolated the bacterium from the CSF of patients with bacterial meningitis. He named the bacterium Diplococcus intracellularis meningitidis.
Biotechnology
Components from Neisseria meningitidis are being harnessed in biotechnology. Its Cas9 enzyme is a useful tool in CRISPR gene editing because the enzyme is small and has distinct targeting features to the commonly used enzyme from Streptococcus pyogenes. The cell-surface protein FrpC from Neisseria meningitidis has been engineered to allow covalent coupling between proteins, because it generates a reactive anhydride when exposed to calcium. The bacterium also expresses unique enzymes able to cleave IgA antibodies.
| Biology and health sciences | Gram-negative bacteria | Plants |
1967137 | https://en.wikipedia.org/wiki/QCD%20vacuum | QCD vacuum | The QCD vacuum is the quantum vacuum state of quantum chromodynamics (QCD). It is an example of a non-perturbative vacuum state, characterized by non-vanishing condensates such as the gluon condensate and the quark condensate in the complete theory which includes quarks. The presence of these condensates characterizes the confined phase of quark matter.
Symmetries and symmetry breaking
Symmetries of the QCD Lagrangian
Like any relativistic quantum field theory, QCD enjoys Poincaré symmetry including the discrete symmetries CPT (each of which is realized). Apart from these space-time symmetries, it also has internal symmetries. Since QCD is an SU(3) gauge theory, it has local SU(3) gauge symmetry.
Since it has many flavours of quarks, it has approximate flavour and chiral symmetry. This approximation is said to involve the chiral limit of QCD. Of these chiral symmetries, the baryon number symmetry is exact. Some of the broken symmetries include the axial U(1) symmetry of the flavour group. This is broken by the chiral anomaly. The presence of instantons implied by this anomaly also breaks CP symmetry.
In summary, the QCD Lagrangian has the following symmetries:
Poincaré symmetry and CPT invariance
SU(3) local gauge symmetry
approximate global SU() × SU() flavour chiral symmetry and the U(1) baryon number symmetry
The following classical symmetries are broken in the QCD Lagrangian:
scale, i.e., conformal symmetry (through the scale anomaly), giving rise to asymptotic freedom
the axial part of the U(1) flavour chiral symmetry (through the chiral anomaly), giving rise to the strong CP problem.
Spontaneous symmetry breaking
When the Hamiltonian of a system (or the Lagrangian) has a certain symmetry, but the vacuum does not, then one says that spontaneous symmetry breaking (SSB) has taken place.
A familiar example of SSB is in ferromagnetic materials. Microscopically, the material consists of atoms with a non-vanishing spin, each of which acts like a tiny bar magnet, i.e., a magnetic dipole. The Hamiltonian of the material, describing the interaction of neighbouring dipoles, is invariant under rotations. At high temperature, there is no magnetization of a large sample of the material. Then one says that the symmetry of the Hamiltonian is realized by the system. However, at low temperature, there could be an overall magnetization. This magnetization has a preferred direction, since one can tell the north magnetic pole of the sample from the south magnetic pole. In this case, there is spontaneous symmetry breaking of the rotational symmetry of the Hamiltonian.
When a continuous symmetry is spontaneously broken, massless bosons appear, corresponding to the remaining symmetry. This is called the Goldstone phenomenon and the bosons are called Goldstone bosons.
Symmetries of the QCD vacuum
The SU() × SU() chiral flavour symmetry of the QCD Lagrangian is broken in the vacuum state of the theory. The symmetry of the vacuum state is the diagonal SU() part of the chiral group. The diagnostic for this is the formation of a non-vanishing chiral condensate , where is the quark field operator, and the flavour index is summed. The Goldstone bosons of the symmetry breaking are the pseudoscalar mesons.
When , i.e., only the up and down quarks are treated as massless, the three pions are the Goldstone bosons. When the strange quark is also treated as massless, i.e., , all eight pseudoscalar mesons of the quark model become Goldstone bosons. The actual masses of these mesons are obtained in chiral perturbation theory through an expansion in the (small) actual masses of the quarks.
In other phases of quark matter the full chiral flavour symmetry may be recovered, or broken in completely different ways.
Experimental evidence
The evidence for QCD condensates comes from two eras, the pre-QCD era 1950–1973 and the post-QCD era, after 1974. The pre-QCD results established that the strong interactions vacuum contains a quark chiral condensate, while the post-QCD results established that the vacuum also contains a gluon condensate.
Motivating results
Gradient coupling
In the 1950s, there were many attempts to produce a field theory to describe the interactions of pions () and nucleons (). The obvious renormalizable interaction between the two objects is the Yukawa coupling to a pseudoscalar:
And this is theoretically correct, since it is leading order and it takes all the symmetries into account. But it doesn't match experiment in isolation. When the nonrelativistic limit of this coupling is taken, the gradient-coupling model is obtained. To lowest order, the nonrelativistic pion field interacts by derivatives. This is not obvious in the relativistic form. A gradient interaction has a very different dependence on the energy of the pion—it vanishes at zero momentum.
This type of coupling means that a coherent state of low momentum pions barely interacts at all. This is a manifestation of an approximate symmetry, a shift symmetry of the pion field. The replacement
leaves the gradient coupling alone, but not the pseudoscalar coupling, at least not by itself. The way nature fixes this in the pseudoscalar model is by simultaneous rotation of the proton-neutron and shift of the pion field. This, when the proper axial SU(2) symmetry is included, is the Gell-Mann Levy σ-model, discussed below.
The modern explanation for the shift symmetry is now understood to be the Nambu-Goldstone non-linear symmetry realization mode, due to Yoichiro Nambu and Jeffrey Goldstone.
The pion field is a Goldstone boson, while the shift symmetry is a manifestation of a degenerate vacuum.
Goldberger–Treiman relation
There is a surprising relationship between the strong interaction coupling of the pions to the nucleons, the coefficient in the nucleon-pion-gradient coupling model, and the axial vector current coefficient of the nucleon, which determines the weak decay rate of the neutron. The relation is
and it is obeyed to 2.5% accuracy.
The constant is the coefficient that determines the neutron decay rate: It gives the normalization of the weak interaction matrix elements for the nucleon. On the other hand, the pion-nucleon coupling is a phenomenological constant describing the (strong) scattering of bound states of quarks and gluons.
The weak interactions are current-current interactions ultimately because they come from a non-Abelian gauge theory. The Goldberger–Treiman relation suggests that the pions, by dint of chiral symmetry breaking, interact as surrogates of sorts of the axial weak currents.
Partially conserved axial current
The structure which gives rise to the Goldberger–Treiman relation was called the partially conserved axial current (PCAC) hypothesis, spelled out in the pioneering σ-model paper. Partially conserved describes the modification of a spontaneously-broken symmetry current by an explicit breaking correction preventing its conservation. The axial current in question is also often called the chiral symmetry current.
The basic idea of SSB is that the symmetry current which performs axial rotations on the fundamental fields does not preserve the vacuum: This means that the current applied to the vacuum produces particles. The particles must be spinless, otherwise the vacuum wouldn't be Lorentz invariant. By index matching, the matrix element must be
where is the momentum carried by the created pion.
When the divergence of the axial current operator is zero, we must have
Hence these pions are massless, , in accordance with Goldstone's theorem.
If the scattering matrix element is considered, we have
Up to a momentum factor, which is the gradient in the coupling, it takes the same form as the axial current turning a neutron into a proton in the current-current form of the weak interaction.
But if a small explicit breaking of the chiral symmetry (due to quark masses) is introduced, as in real life, the above divergence does not vanish, and the right hand side involves the mass of the pion, now a Pseudo-Goldstone boson.
Soft pion emission
Extensions of the PCAC ideas allowed Steven Weinberg to calculate the amplitudes for collisions which emit low energy pions from the amplitude for the same process with no pions. The amplitudes are those given by acting with symmetry currents on the external particles of the collision.
These successes established the basic properties of the strong interaction vacuum well before QCD.
Pseudo-Goldstone bosons
Experimentally it is seen that the masses of the octet of pseudoscalar mesons is very much lighter than the next lightest states; i.e., the octet of vector mesons (such as the rho meson). The most convincing evidence for SSB of the chiral flavour symmetry of QCD is the appearance of these pseudo-Goldstone bosons. These would have been strictly massless in the chiral limit. There is convincing demonstration that the observed masses are compatible with chiral perturbation theory. The internal consistency of this argument is further checked by lattice QCD computations which allow one to vary the quark mass and check that the variation of the pseudoscalar masses with the quark mass is as required by chiral perturbation theory.
Eta prime meson
This pattern of SSB solves one of the earlier "mysteries" of the quark model, where all the pseudoscalar mesons should have been of nearly the same mass. Since , there should have been nine of these. However, one (the SU(3) singlet η′ meson) has quite a larger mass than the SU(3) octet. In the quark model, this has no natural explanation – a mystery named the η−η′ mass splitting (the η is one member of the octet, which should have been degenerate in mass with the η′).
In QCD, one realizes that the η′ is associated with the axial UA(1) which is explicitly broken through the chiral anomaly, and thus its mass is not "protected" to be small, like that of the η. The η–η′ mass splitting can be explained
through the 't Hooft instanton mechanism,
whose realization is also known as Witten–Veneziano mechanism.
Current algebra and QCD sum rules
PCAC and current algebra also provide evidence for this pattern of SSB. Direct estimates of the chiral condensate also come from such analysis.
Another method of analysis of correlation functions in QCD is through an operator product expansion (OPE). This writes the vacuum expectation value of a non-local operator as a sum over VEVs of local operators, i.e., condensates. The value of the correlation function then dictates the values of the condensates. Analysis of many separate correlation functions gives consistent results for several condensates, including the gluon condensate, the quark condensate, and many mixed and higher order condensates. In particular one obtains
Here refers to the gluon field tensor, to the quark field, and to the QCD coupling.
These analyses are being refined further through improved sum rule estimates and direct estimates in lattice QCD. They provide the raw data which must be explained by models of the QCD vacuum.
Models
A full solution of QCD should give a full description of the vacuum, confinement and the hadron spectrum. Lattice QCD is making rapid progress towards providing the solution as a systematically improvable numerical computation. However, approximate models of the QCD vacuum remain useful in more restricted domains. The purpose of these models is to make quantitative sense of some set of condensates and hadron properties such as masses and form factors.
This section is devoted to models. Opposed to these are systematically improvable computational procedures such as large QCD and lattice QCD, which are described in their own articles.
The Savvidy vacuum, instabilities and structure
The Savvidy vacuum is a model of the QCD vacuum which at a basic level is a statement that it cannot be the conventional Fock vacuum empty of particles and fields. In 1977, George Savvidy showed that the QCD vacuum with zero field strength is unstable, and decays into a state with a calculable non vanishing value of the field. Since condensates are scalar, it seems like a good first approximation that the vacuum contains some non-zero but homogeneous field which gives rise to these condensates. However, Stanley Mandelstam showed that a homogeneous vacuum field is also unstable. The instability of a homogeneous gluon field was argued by Niels Kjær Nielsen and Poul Olesen in their 1978 paper. These arguments suggest that the scalar condensates are an effective long-distance description of the vacuum, and at short distances, below the QCD scale, the vacuum may have structure.
The dual superconducting model
In a type II superconductor, electric charges condense into Cooper pairs. As a result, magnetic flux is squeezed into tubes. In the dual superconductor picture of the QCD vacuum, chromomagnetic monopoles condense into dual Cooper pairs, causing chromoelectric flux to be squeezed into tubes. As a result, confinement and the string picture of hadrons follows. This dual superconductor picture is due to Gerard 't Hooft and Stanley Mandelstam. 't Hooft showed further that an Abelian projection of a non-Abelian gauge theory contains magnetic monopoles.
While the vortices in a type II superconductor are neatly arranged into a hexagonal or occasionally square lattice, as is reviewed in Olesen's 1980 seminar one may expect a much more complicated and possibly dynamical structure in QCD. For example, nonabelian Abrikosov-Nielsen-Olesen vortices may vibrate wildly or be knotted.
String models
String models of confinement and hadrons have a long history. They were first invented to explain certain aspects of crossing symmetry in the scattering of two mesons. They were also found to be useful in the description of certain properties of the Regge trajectory of the hadrons. These early developments took on a life of their own called the dual resonance model (later renamed string theory). However, even after the development of QCD string models continued to play a role in the physics of strong interactions. These models are called non-fundamental strings or QCD strings, since they should be derived from QCD, as they are, in certain approximations such as the strong coupling limit of lattice QCD.
The model states that the colour electric flux between a quark and an antiquark collapses into a string, rather than spreading out into a Coulomb field as the normal electric flux does. This string also obeys a different force law. It behaves as if the string had constant tension, so that separating out the ends (quarks) would give a potential energy increasing linearly with the separation. When the energy is higher than that of a meson, the string breaks and the two new ends become a quark-antiquark pair, thus describing the creation of a meson. Thus confinement is incorporated naturally into the model.
In the form of the Lund model Monte Carlo program, this picture has had remarkable success in explaining experimental data collected in electron-electron and hadron-hadron collisions.
Bag models
Strictly, these models are not models of the QCD vacuum, but of physical single particle quantum states — the hadrons. The model proposed originally in 1974 by A. Chodos et al.
consists of inserting a quark model in a perturbative vacuum inside a volume of space called a bag. Outside this bag is the real QCD vacuum, whose effect is taken into account through the difference between energy density of the true QCD vacuum and the perturbative vacuum (bag constant ) and boundary conditions imposed on the quark wave functions and the gluon field. The hadron spectrum is obtained by solving the Dirac equation for quarks and the Yang–Mills equations for gluons. The wave functions of the quarks satisfy the boundary conditions of a fermion in an infinitely deep potential well of scalar type with respect to the Lorentz group. The boundary conditions for the gluon field are those of the dual color superconductor. The role of such a superconductor is attributed to the physical vacuum of QCD. Bag models strictly prohibit the existence of open color (free quarks, free gluons, etc.) and lead in particular to string models of hadrons.
The chiral bag model couples the axial vector current of the quarks at the bag boundary to a pionic field outside of the bag. In the most common formulation, the chiral bag model basically replaces the interior of the skyrmion with the bag of quarks. Very curiously, most physical properties of the nucleon become mostly insensitive to the bag radius. Prototypically, the baryon number of the chiral bag remains an integer, independent of bag radius: the exterior baryon number is identified with the topological winding number density of the Skyrme soliton, while the interior baryon number consists of the valence quarks (totaling to one) plus the spectral asymmetry of the quark eigenstates in the bag. The spectral asymmetry is just the vacuum expectation value summed over all of the quark eigenstates in the bag. Other values, such as the total mass and the axial coupling constant , are not precisely invariant like the baryon number, but are mostly insensitive to the bag radius, as long as the bag radius is kept below the nucleon diameter. Because the quarks are treated as free quarks inside the bag, the radius-independence in a sense validates the idea of asymptotic freedom.
Instanton ensemble
Another view states that BPST-like instantons play an important role in the vacuum structure of QCD. These instantons were discovered in 1975 by Alexander Belavin, Alexander Markovich Polyakov, Albert S. Schwarz and Yu. S. Tyupkin as topologically stable solutions to the Yang-Mills field equations. They represent tunneling transitions from one vacuum state to another. These instantons are indeed found in lattice calculations. The first computations performed with instantons used the dilute gas approximation. The results obtained did not solve the infrared problem of QCD, making many physicists turn away from instanton physics. Later, though, an instanton liquid model was proposed, turning out to be more promising an approach.
The dilute instanton gas model departs from the supposition that the QCD vacuum consists of a gas of BPST-like instantons. Although only the solutions with one or few instantons (or anti-instantons) are known exactly, a dilute gas of instantons and anti-instantons can be approximated by considering a superposition of one-instanton solutions at great distances from one another. Gerard 't Hooft calculated the effective action for such an ensemble, and he found an infrared divergence for big instantons, meaning that an infinite amount of infinitely big instantons would populate the vacuum.
Later, an instanton liquid model was studied. This model starts from the assumption that an ensemble of instantons cannot be described by a mere sum of separate instantons. Various models have been proposed, introducing interactions between instantons or using variational methods (like the "valley approximation") endeavoring to approximate the exact multi-instanton solution as closely as possible. Many phenomenological successes have been reached. Whether an instanton liquid can explain confinement in 3+1 dimensional QCD is not known, but many physicists think that it is unlikely.
Center vortex picture
A more recent picture of the QCD vacuum is one in which center vortices play an important role. These vortices are topological defects carrying a center element as charge. These vortices are usually studied using lattice simulations, and it has been found that the behavior of the vortices is closely linked with the confinement–deconfinement phase transition: in the confinement phase vortices percolate and fill the spacetime volume, in the deconfinement phase they are much suppressed. Also it has been shown that the string tension vanished upon removal of center vortices from the simulations, hinting at an important role for center vortices.
| Physical sciences | Quantum mechanics | Physics |
8528731 | https://en.wikipedia.org/wiki/Carassius | Carassius | Carassius is a genus in the ray-finned fish family Cyprinidae. Most species in this genus are commonly known as crucian carps, though that term often refers specifically to C. carassius. The most well known species is the goldfish (C. auratus). They have a Eurasian distribution, apparently originating further to the west than the typical carps (Cyprinus genus, which includes the common carp).
Species of Carassius genus are not closely related of the typical carps of Cyprinus genus, but rather form a more basal lineage of the subfamily Cyprininae.
Species
Carassius contains the following species:
Carassius auratus (Linnaeus, 1758) (Goldfish)
Carassius carassius (Linnaeus, 1758) (Crucian carp)
Carassius cuvieri Temminck & Schlegel, 1846 (Japanese white crucian carp)
Carassius gibelio (Bloch, 1782) (Prussian carp)
Carassius langsdorfii Temminck & Schlegel, 1846 (Ginbuna)
Carassius praecipuus Kottelat, 2017
| Biology and health sciences | Cypriniformes | Animals |
6567126 | https://en.wikipedia.org/wiki/Cubic%20metre%20per%20second | Cubic metre per second | Cubic metre per second or cubic meter per second in American English (symbol m3s−1 or m3/s) is the unit of volumetric flow rate in the International System of Units (SI). It corresponds to the exchange or movement of the volume of a cube with sides of in length (a cubic meter, originally a stere) each second. It is popularly used for water flow, especially in rivers and streams, and fractions for HVAC values measuring air flow.
The term cumec is sometimes used as an acronym for full unit name, with the plural form cumecs also common in speech. It is commonly used between workers in the measurement of water flow through natural streams and civil works, but rarely used in writing.
Data in units of m3s−1 are used along the y-axis or vertical axis of a flow hydrograph, which describes the time variation of discharge of a river (the mean velocity multiplied by cross-sectional area). A moderately sized river discharges in the order of 100 m3s−1.
Conversions
| Physical sciences | Volumetric flow rate | Basics and measurement |
20063724 | https://en.wikipedia.org/wiki/Colored%20gold | Colored gold | Colored gold is the name given to any gold that has been treated using techniques to change its natural color. Pure gold is slightly reddish yellow in color, but colored gold can come in a variety of different colors by alloying it with different elements.
Colored golds can be classified in three groups:
Alloys with silver and copper in various proportions, producing white, yellow, green and red golds. These are typically malleable alloys.
Intermetallic compounds, producing blue and purple golds, as well as other colors. These are typically brittle, but can be used as gems and inlays.
Surface treatments, such as oxide layers.
Pure 100% (in practice, 99.9% or better) gold is 24 karat by definition, so all colored golds are less pure than this, commonly 18K (75%), 14K (58.5%), 10K (41.6%), or 9K (37.5%).
Alloys
White gold
The word white covers a broad range of colors that borders or overlaps pale yellow, tinted brown, and even very pale rose. White gold is an alloy of gold and at least one white metal (usually nickel, silver, platinum or palladium). Like yellow gold, the purity of white gold is given in karats.
White gold's properties vary depending on the metals used and their proportions. A common white gold formulation consists of 90% wt. gold and 10% wt. nickel. Copper can be added to increase malleability. To give it a shiny finish and brightness sometimes it plated with rhodium. The alloys used in the jewelry industry are gold-palladium-silver and gold-nickel-copper-zinc. Palladium and nickel act as primary bleaching agents for gold; zinc acts as a secondary bleaching agent to attenuate the color of copper. As a result, white gold alloys can be used for many different purposes. Nickel alloys are hard and strong, and therefore good for rings and pins. Gold-palladium alloys are soft, pliable, and good for white-gold gemstone settings. The strength of gold-nickel-copper alloys is caused by formation of two phases: a gold-rich Au-Cu, and a nickel-rich Ni-Cu, and the resulting hardening of the material.
The nickel used in some white gold alloys can cause an allergic reaction when worn over long periods (also notably on some wristwatch casings). This reaction, typically a minor skin rash from nickel dermatitis, occurs in about one out of eight people; because of this, many countries do not use nickel in their white gold formulations.
Rose, red, and pink gold
Rose gold is a gold-copper alloy widely used for specialized jewelry. Rose gold, also known as pink gold and red gold, was popular in Russia at the beginning of the 19th century, and was also known as Russian gold. Rose gold jewelry is becoming more popular in the 21st century, and is commonly used for wedding rings, bracelets, and other jewelry.
Although the names are often used interchangeably, the difference between red, rose, and pink gold is the copper content: the higher the copper content, the stronger the red coloration. Pink gold uses the least copper, followed by rose gold, with red gold having the highest copper content. Examples of the common alloys for 18K rose gold, 18K red gold, 18K pink gold, and 12K red gold include:
18K red gold: 75% gold, 25% copper
18K rose gold: 75% gold, 22.25% copper, 2.75% silver
18K pink gold: 75% gold, 20% copper, 5% silver
12K red gold: 50% gold and 50% copper
Up to 15% zinc can be added to copper-rich alloys to change their color to reddish yellow or dark yellow. 14K red gold, often found in the Middle East, contains 41.67% copper.
The highest karat version of rose gold, also known as crown gold, is 22 karat.
Amongst the alloys made of gold, silver, and copper, the hardest is the 18.1 K pink gold (75.7% gold and 24.3% copper). An alloy with only gold and silver is the hardest at 15.5 K (64.5% gold and 35.5% silver).
During ancient times, due to impurities in the smelting process, gold frequently turned a reddish color. This is why many Greek and Roman texts, and some texts from the Middle Ages, describe gold as "red".
Spangold
Some gold-copper-aluminium alloys form a fine surface texture at heat treatment, yielding a spangling effect. At cooling, they undergo a quasi-martensitic transformation from body-centered cubic to body-centered tetragonal phase; the transformation does not depend on the cooling rate. A polished object is heated in hot oil to 150–200 °C for 10 minutes then cooled below 20 °C, forming a sparkly surface covered with tiny facets.
The alloy of 76% gold, 19% copper, and 5% aluminium yields a yellow color; the alloy of 76% gold, 18% copper, and 6% aluminium is pink.
Green gold
Electrum, a naturally occurring alloy of silver and gold, develops a greenish cast with increasing silver content, ranging in color from green-yellow (for proportions of silver between 14% and 29%) to pale green-yellow (for proportions of silver between 29% and 50%). It was known to the ancient Persians as long ago as 860 BC. However, electrum was used even thousands of years before that, by both the Akkadians and Ancient Egyptians (as evidenced by the Royal Cemetery at Ur). Even the tops of some Egyptian pyramids were known to be capped in thin layers of electrum. Fired enamels adhere better to these alloys than to pure gold.
Cadmium can also be added to gold alloys to create a green color, but there are health concerns regarding its use, as cadmium is highly toxic. Adding 2% cadmium to 18K red gold yields a light green color, whereas the alloy of 75% gold, 15% silver, 6% copper, and 4% cadmium is dark green.
Gray gold
Gray gold alloys are usually made from gold and palladium. A cheaper alternative which does not use palladium is made by adding silver, manganese, and copper to the gold in specific ratios.
Intermetallic
All the AuX2 intermetallics have the fluorite (CaF2) crystal structure, and, therefore, are brittle. Deviation from the stoichiometry results in loss of color. Slightly nonstoichiometric compositions are used, however, to achieve a fine-grained two- or three-phase microstructure with reduced brittleness. Another way of reducing brittleness is to add a small amount of palladium, copper, or silver.
The intermetallic compounds tend to have poor corrosion resistance. The less noble elements are leached to the environment, and a gold-rich surface layer is formed. Direct contact of blue and purple gold elements with skin should be avoided as exposure to sweat may result in metal leaching and discoloration of the metal surface.
Purple gold
Purple gold (also called amethyst gold and violet gold) is an alloy of gold and aluminium rich in gold–aluminium intermetallic (AuAl2). Gold content in AuAl2 is around 79% and can therefore be referred to as 18 karat gold. Purple gold is more brittle than other gold alloys (called the "purple plague" when it forms and causes serious faults in electronics), as it is an intermetallic compound instead of a malleable alloy, and a sharp blow may cause it to shatter. It is therefore usually machined and faceted to be used as a "gem" in conventional jewelry rather than by itself. At a lower content of gold, the material is composed of the intermetallic and an aluminium-rich solid solution phase. At a higher content of gold, the gold-richer intermetallic AuAl forms; the purple color is preserved to about 15% of aluminium. At 88% of gold the material is composed of AuAl and changes color. The actual composition of AuAl2 is closer to Au6Al11 as the sublattice is incompletely occupied.
Blue gold
Blue gold is an alloy of gold and either gallium or indium. Gold-indium contains 46% gold (about 11 karat) and 54% indium, forming an intermetallic compound AuIn2. While several sources remark this intermetallic to have "a clear blue color", in fact the effect is slight: AuIn2 has CIE LAB color coordinates of 79, −3.7, −4.2 which appears roughly as a grayish color. With gallium, gold forms an intermetallic AuGa2 (58.5% Au, 14ct) which has slighter bluish hue. The melting point of AuIn2 is 541 °C, for AuGa2 it is 492 °C. AuIn2 is less brittle than AuGa2, which itself is less brittle than AuAl2.
A surface plating of blue gold on karat gold or sterling silver can be achieved by a gold plating of the surface, followed by indium plating, with layer thickness matching the 1:2 atomic ratio. A heat treatment then causes interdiffusion of the metals and formation of the required intermetallic compound.
Surface treatments
Black gold
Black gold is a type of gold used in jewelry. Black-colored gold can be produced by various methods:
Patination by applying sulfur- and oxygen-containing compounds.
Plasma-assisted chemical vapor deposition process involving amorphous carbon
Controlled oxidation of gold containing chromium or cobalt (e.g. 75% gold, 25% cobalt).
A range of colors from brown to black can be achieved on copper-rich alloys by treatment with potassium sulfide.
Cobalt-containing alloys, e.g. 75% gold with 25% cobalt, form a black oxide layer with heat treatment at 700–950 °C. Copper, iron and titanium can be also used for such effect. Gold-cobalt-chromium alloy (75% gold, 15% cobalt, 10% chromium) yields a surface oxide that is olive-tinted because of the chromium(III) oxide content, is about five times thinner than Au-Co and has significantly better wear resistance. The gold-cobalt alloy consists of gold-rich (about 94% Au) and cobalt-rich (about 90% Co) phases; the cobalt-rich phase grains are capable of oxide-layer formation on their surface.
More recently, black gold can be formed by creating nanostructures on the surface. A femtosecond laser pulse deforms the surface of the metal, creating an immensely increased surface area which absorbs virtually all the light that falls on it, thus rendering it deep black, but this method is used in high technology applications rather than for appearance in jewelry. The blackness is due to the excitation of localized surface plasmons which creates strong absorption in a broad range in plasmon resonance. The broadness of the plasmon resonance, and absorption wavelength range, depends on the interaction between different gold nanoparticles.
Blue gold
Oxide layers can also be used to obtain blue gold from an alloy of 75% gold, 24.4% iron, and 0.6% nickel; the layer forms on heat treatment in air between 450 and 600 °C.
A rich sapphire blue colored gold of 20–23K can also be obtained by alloying with ruthenium, rhodium, and three other elements and heat-treating at 1800 °C, to form the 3–6 micrometers thick colored surface oxide layer.
| Physical sciences | Specific alloys | Chemistry |
22641228 | https://en.wikipedia.org/wiki/Transplanter | Transplanter | A transplanter is an agricultural machine used for transplanting seedlings to the field. Transplanters greatly reduce time required to transplant seedlings compared to manual transplanting. Among the crops that are transplanted with transplanters are strawberries, vegetables, tomatoes, cabbages, tobacco and rice.
Semi-automatic mechanical transplanters are a common type, which can be self-propelled, or towed by a tractor at a low speed. A row of three to six operators feed seedlings from germination trays into hoppers which feed into the delivery mechanism.
Gallery
| Technology | Farm and garden machinery | null |
24130500 | https://en.wikipedia.org/wiki/Stoney%20units | Stoney units | In physics, the Stoney units form a system of units named after the Irish physicist George Johnstone Stoney, who first proposed them in 1881. They are the earliest example of natural units, i.e., a coherent set of units of measurement designed so that chosen physical constants fully define and are included in the set.
Units
The constants that Stoney used to define his set of units is the following:
, the speed of light in vacuum,
, the gravitational constant,
, the Coulomb constant,
, the charge on the electron.
Later authors typically replace the Coulomb constant with .
This means that the numerical values of all these constants, when expressed in coherent Stoney units, is equal to one:
In Stoney units, the numerical value of the reduced Planck constant is
where is the fine-structure constant.
History
George Stoney was one of the first scientists to understand that electric charge was quantized; from this quantization and three other constants that he perceived as being universal (a speed from electromagnetism, and the coefficients in the electrostatic and gravitational force equations) he derived the units that are now named after him.
Stoney's derived estimate of the unit of charge, 10−20 ampere-second, was of the modern value of the charge of the electron due to Stoney using the approximated value of 1018 for the number of molecules presented in one cubic millimetre of gas at standard temperature and pressure. Using the modern values for the Avogadro constant and for the volume of a gram-molecule under these conditions of , the modern value is , instead of Stoney's 1018.
Stoney units and Planck units
Stoney's set of base units is similar to the one used in Planck units, proposed independently by Planck thirty years later, in which Planck normalized the Planck constant in place of the elementary charge.
Planck units are more commonly used than Stoney units in modern physics, especially for quantum gravity (including string theory). Rarely, Planck units are referred to as Planck–Stoney units.
The Stoney length and the Stoney energy, collectively called the Stoney scale, are not far from the Planck length and the Planck energy, the Planck scale. The Stoney scale and the Planck scale are the length and energy scales at which quantum processes and gravity occur together. At these scales, a unified theory of physics is thus required. The only notable attempt to construct such a theory from the Stoney scale was that of Hermann Weyl, who associated a gravitational unit of charge with the Stoney length
and who appears to have inspired Dirac's fascination with the large numbers hypothesis.
Since then, the Stoney scale has been largely neglected in the development of modern physics, although it is still occasionally discussed.
The ratio of Stoney units to Planck units of length, time and mass is , where is the fine-structure constant:
| Physical sciences | Measurement systems | Basics and measurement |
25532030 | https://en.wikipedia.org/wiki/Nonlinear%20metamaterial | Nonlinear metamaterial | A nonlinear metamaterial is an artificially constructed material that can exhibit properties not yet found in nature. Its response to electromagnetic radiation can be characterized by its permittivity and material permeability. The product of the permittivity and permeability results in the refractive index. Unlike natural materials, nonlinear metamaterials can produce a negative refractive index. These can also produce a more pronounced nonlinear response than naturally occurring materials.
Nonlinear metamaterials are a periodic, nonlinear, transmission medium. These are a type of negative index metamaterial where the nonlinearity is available because the microscopic electric field of the inclusions can be larger than the macroscopic electric field of the electromagnetic (EM) source. This then becomes a useful tool which allows for enhancing the nonlinear behavior of the metamaterial. A dominant nonlinear response, however, can be derived from the hysteresis-type dependence of the material's magnetic permeability on the magnetic component of the incident electromagnetic wave (light) propagating through the material. Furthermore, the hysteresis-type dependence of the magnetic permeability on the field intensity allows changing the material from left to right-handed and back.
Nonlinear media are essential for nonlinear optics. However most optical materials have a relatively weak nonlinear response, meaning that their properties only change by a small amount for large changes in intensity of the electromagnetic field. Nonlinear metamaterials can overcome this limitation, since the local fields of the resonant structures can be much larger than the average value of the field - in this respect metamaterials are similar to other composite media, such e.g. as random metal-dielectric composites, including fractal clusters and semicoutinouos/percolation metal films, where the areas with enhanced local light fields - “hot spots” - produce giant linear and non-linear optical responses.
Overview of metamaterials
Metamaterials are incarnations of materials first proposed by a Russian theorist, Victor Veselago in 1967.
Nonlinear metamaterials, a type of metamaterial, are being developed in order to manipulate electromagnetic radiation in new ways. Optical and electromagnetic properties of natural materials are often altered through chemistry. With metamaterials optical and electromagnetic properties can be engineered through the geometry of its unit cells. The unit cells are materials that are ordered in geometric arrangements with dimensions that are fractions of the wavelength of the radiated electromagnetic wave.
By having the freedom to alter effects by adjusting the configurations and sizes of the unit cells, control over permittivity and magnetic permeability can be achieved. These two parameters (or quantities) determine the propagation of electromagnetic waves in matter. Therefore, the achievable electromagnetic and optical effects can be extended.
Optical properties can be expanded beyond the capabilities of lenses, mirrors, and other conventional materials. One of the effects most studied is the negative index of refraction first proposed by Victor Veselago in 1967. Negative index materials, exhibit optical properties opposite to those of glass, air, and the other conventional materials. At the correct frequencies, the negative index materials refract electromagnetic waves in novel ways, to a zero index or negative index. Also, energy can propagate in the opposite direction which can result in compensation mechanisms, among other possibilities.
Interactions
Materials which scatter light or other electromagnetic waves create a general physical process where the different frequencies of light are forced to deviate from a straight trajectory. It is because, physically, the material is non-uniform at one, or more, or many places.
Furthermore, the optical sciences make predictions about the path of light traversing through a material. When light deviates from its predicted (reflected) path, this also is considered scattering. The split ring resonators which make up metamaterials are engineered to scatter light at resonance. Moreover, these resonant scattering elements are purposely designed at a uniform size throughout the material. This uniform size is much smaller than the wavelength of the frequency of light propagating through the material.
Since the repeating, scattering, resonant elements, which make up the engineered material are much smaller than the frequency of propagating light, metamaterials can now, also, be described in terms of macroscopic quantities. This description is simply another way to view metamaterials. And these are electric permittivity, ε and magnetic permeability, μ.
Hence, by designing the individual, geometrically shaped unit of the material, called a cell, as the right kind of composite, it becomes a material with macroscopic properties that do not occur in nature.
Of particular interest regarding nonlinear metamaterials, is the artificially induced macroscopic property known as negative refractive index. This effect is created by Negative index metamaterials (NIMs), which are employed for use as nonlinear metamaterials. In NIMs, nonlinear phenomena such as second-harmonic generation and parametric amplification can take on highly unusual characteristics. Namely, the fact that the wavevector and the Poynting vector of a wave propagating in a NIM are counter-directed alters the phase-matching conditions for the interacting waves, resulting in backward propagating waves as well as considerably changed Manley-Rowe relations and the distribution in space of the interacting fields' intensity.
Non-linear properties of left-handed metamaterials
Previous studies of left-handed or negative index metamaterials were focused on the linear properties of the medium during wave propagation. In such cases, the view was that magnetic permeability and material permittivity are each not dependent on the intensity of the electromagnetic field. However, creating tunable structures requires knowledge of non-linear properties where the intensity of the electromagnetic field alters the permittivity, or permeability, or both, which in turn affects the range of transmission spectra or stop band spectra. Hence, the effective permeability is dependent on the macroscopic magnetic field intensity. As the field intensity is varied, switching between its positive and negative values can occur. Consequently, the material can switch from being left-handed to being right-handed, or vice versa.
A composite structure consisting of a square lattice of the periodic arrays of conducting wires and split-ring resonators, produces an enhanced magnetic response. Without the correct magnetic response, it is not possible to produce a left-handed material.
Tunable split-ring resonators for nonlinear negative-index metamaterials
Variable capacitance diodes are incorporated into the split-ring cell producing a dynamic tunable system.
Reconfigurable refractive index (infrared)
Source radiation of near infrared wavelengths are applied to a metamaterial system. The index of refraction can be reconfigured to exhibit negative values, zero, or positive values.
SRR microwave nonlinear tunable metamaterial
Fabrication and experimental studies of the properties of the first nonlinear tunable metamaterial were operating at microwave frequencies. Such a metamaterial was fabricated by modifying the properties of SRRs and introducing varactor diodes in each SRR element of the composite structure such that the whole structure becomes dynamically tunable by varying the amplitude of the propagating electromagnetic waves. In particular, the power dependent transmission of the left-handed and magnetic metamaterials at higher powers were demonstrated, as well as the generation of particular harmonics, as was theoretically suggested earlier.
SRR microwave nonlinear magnetic metamaterials
The fabrication and experimental studies of the properties of thenonlinear tunable magnetic metamaterial were operating at microwave frequencies. Varactor diodes are symmetrically introduced, which results in dynamic tunability for the whole structure. Since the magnetic component of the interaction determines the application, the power dependency is demonstrated. Nonlinearity-dependent enhancement or suppression of the transmission turns out to be dynamically tunable.
SRR microwave nonlinear electric metamaterials
A novel class of nonlinear metamaterials has been proposed and engineered to demonstrate a resonant electric response within microwave frequency ranges. These metamaterials incorporate varactor diodes as nonlinear components within each resonator. This design enables the manipulation of the frequency of the electric mode stop band by modulating the incident power levels. Importantly, this modulation does not impact the magnetic response characteristics of the metamaterial. These elements could be combined with the previously developed nonlinear magnetic metamaterials in order to create negative index media with a control over both electric and magnetic nonlinearities.
Nonlinear resonators are designed in a similar fashion. A strong nonlinear electric response is obtained.
Sub-diffraction limit for non-linear metamaterial lens
By covering a thin flat nonlinear lens on the sources, the sub-diffraction-limit observation can be achieved by measuring either the near-field distribution or the far-field radiation of the sources at the harmonic frequencies and calculating the IFT to obtain the sub-wavelength imaging. The higher order harmonics are used, the higher resolution is obtained.
Non-linear electric metamaterial
A new type of nonlinear metamaterial is designed, and analyzed with a dominant negative electric response. Introducing nonlinearity into the electric response makes it tunable while leaving the magnetic response unchanged. A nonlinear NIM containing tunable electric and magnetic elements, which can respond independently is possible.
EM field shielding by non-linear metamaterials
It is well known that over certain frequencies, typical metals can reflect electromagnetic
(EM) fields and can thus be used as electromagnetic shielding materials. However, conventional linear LHMs cannot be used to shield electromagnetic fields. This is drastically modified when nonlinearity of the magnetic response is taken into account, creating a controllable shielding effect in LHMs, accompanied by a parametric reflection.
Meta-dimer metamaterial
A meta-dimer is composed of two spatially separated SRRs, with the two SRRs identical in each unit cell. The proximity of the SRRs in the dimer results in relatively strong coupling between them. A metamaterial comprising a large number of such metadimers can be utilized as an actively tunable medium at optical wavelengths. If either or both of the SRRs in the meta-dimer become nonlinear, the metamaterial itself acquires nonlinear properties. This can allow for nonlinear behavior, such as tunability in real time. Stereometamaterials are also a type of meta-dimer.
| Physical sciences | Basics_3 | Physics |
28661667 | https://en.wikipedia.org/wiki/Parmotrema%20perlatum | Parmotrema perlatum | Parmotrema perlatum, commonly known as the powdered ruffle lichen, is a common species of foliose lichen in the family Parmeliaceae. The species has a cosmopolitan distribution and occurs throughout the Northern and Southern Hemispheres. Parmotrema perlatum is a prominent and widely recognised species within its genus across primarily temperate zones, preferring humid, oceanic-suboceanic habitats. It is found in diverse geographic areas including Africa, North and South America, Asia, Australasia, Europe, and islands in the Atlantic and Pacific oceans. It usually grows on bark, but occasionally occurs on siliceous rocks, often among mosses.
The thallus of Parmotrema perlatum is large, light-grey to pale-blue patch-shaped with rounded and ruffled and often with black hair- at the edges. Distinguishing features of the lichen include its conspicuous soralia (reproductive structures) near the lobe edges, curled leaf-like lobes, and a narrow, shiny, and sometimes wrinkly area on the underside near the margin. This species is known for producing certain secondary metabolites, namely atranorin and a group of substances known as the stictic acid complex, which includes stictic and constictic acids, among other related compounds. These morphological and chemical characteristics help distinguish P. perlatum from several other potential lookalikes.
Parmotrema perlatum has a complex taxonomic history, having undergone multiple reclassifications since its original description in 1762. Significant efforts in the mid-20th century helped clarify its nomenclature, stabilising its current name. Although there were challenges to this name in the 1980s, it was confirmed as valid in 2004. More recently, DNA studies suggest that there may be hidden diversity within the species, indicating the need for further taxonomic evaluation.
The lichen is used as a spice in Indian cuisine. For this purpose, it is commonly known as black stone flower or kalpasi (among other names). Although nearly tasteless on its own, it releases an earthy fragrance and taste when cooked in with oil or butter.
Systematics
Historical taxonomy
The taxonomy of Parmotrema perlatum has a rich history marked by periods of confusion and clarification that typify the dynamic nature of botanical classification. It was originally described as Lichen perlatus by William Hudson in his 1762 work Flora Anglica. Hudson described it as a foliaceous (leafy) lichen with creeping, lobed, and smooth characteristics, having a pearly edge, a farinaceous () texture, and a black underside, adorned with slightly scalloped, brown, stalked fruiting bodies. The taxon was later transferred to the genus Parmelia by Erik Acharius in 1803, becoming Parmelia perlata. The name was well-established in scientific literature, being cited extensively in works like Alexander Zahlbruckner's popular 1929 catalogue.
In 1952, Maurice Choisy reclassified it under the current name, Parmotrema perlatum. The nomenclature of Parmotrema perlatum was revisited in the late 20th century, amid a broader effort to clarify the typification and application of early lichen names. Mason Hale, in 1961, undertook a detailed restudy of the species, selecting a lectotype from the Dillenian collections—the herbarium and associated works of Johann Jacob Dillenius housed at the University of Oxford. This solidified the application of Hudson's name and was part of a larger trend in lichenology to fix historical names to specific herbarium specimens to stabilise nomenclature.
The name Parmelia perlata was widely accepted until Hale and Ahti (1986) encountered the designation Lichen chinensis, introduced by Pehr Osbeck in 1757. They proposed the name Parmotrema chinense, based on the assumption that Osbeck's specimen corresponded to the well-known species Parmotrema perlatum. However, this proposal was not universally adopted due to the lack of valid typification and the name's absence in the literature between 1757 and 1986.
Impact of the Tokyo Code (1993)
The Tokyo Code of 1993 extended the provisions for conserving names to all species, not just those of major economic importance. This change in the International Code of Botanical Nomenclature allowed for the conservation of names that would promote nomenclatural stability. Despite this provision, no formal proposal was made to conserve the name Parmotrema chinense, and thus it did not gain widespread acceptance.
David Hawksworth's 2004 study brought significant clarity to the taxonomic confusion. He rediscovered Osbeck's original material in Linnaeus' herbarium and identified it as belonging to Parmotrema tinctorum, not Parmotrema perlatum. Hawksworth demonstrated that Lichen chinensis was not validly published because it lacked a proper description and was linked with an expression of doubt by Osbeck. Hawksworth's work led to the reinstatement of the name Parmotrema perlatum, confirming that Hudson's name was legitimate and should continue to be used. This resolution was based on the original typification by Hale and the invalid publication status of Lichen chinensis.
Recent studies suggest that the circumscription of Parmotrema perlatum may need to be revised. Research utilising DNA sequencing has uncovered cryptic diversity within the genus Parmotrema, indicating that traditional phenotype-based identification methods may underestimate species diversity. Specifically, the genetic analysis of P. perlatum and related species revealed multiple distinct lineages that were previously grouped under a single nominal taxon. These findings highlight the need for a comprehensive taxonomic re-evaluation of P. perlatum to accurately delineate species boundaries and account for hidden genetic diversity.
Phylogeny
In molecular phylogenetics analysis, Parmotrema perlatum has a sister relationship with Parmotrema crinitum. These two species form a clade that itself is sister to a clade with P. austrosinense and P. tinctorum. In a comprehensive phylogenetic analysis by Stelate and colleagues (2022), P. perlatum and P crinitum were found to form a well-supported monophyletic group using internal transcribed spacer sequences and several analytical methods. This study highlights the complexity of species boundaries within the genus and the need for further research incorporating additional molecular markers to confirm these findings.
Common names
Vernacular names used for this species include black stone flower, stone lichen, sea lichen, kalpasi, kalpas, kalpashi, and kalpash. The latter name and its variations, however, have been used as a crude drug in Indian medicines for more than one species, including Parmotrema perlatum, Parmotrema tinctorum, and Everniastrum cirrhatum. In North America, vernacular names used for the species include "powdered ruffle lichen", "powdered scatter-rug", and "queen ruffle". The species epithet perlatum refers to the pearl-like margins of the lobes, which are directly referenced in Hudson's original 1762 description of the species. He proposed the English name "pearl lichen"; this name later morphed into "pearly parmelia" in some 19th-century British accounts of lichen flora.
Description
Parmotrema perlatum has a thallus that ranges from loosely to tightly attached to the surface it grows on, forming expansive, spreading colonies that often merge together. Individual thalli typically measure up to in diameter. The upper thallus surface is greenish-grey, blue-grey, or yellowish-grey in colour, lacking and either free of spots (), or with few maculae. This species develops soredia, a type of asexual reproductive structure, aiding in its propagation. The of this lichen vary from 1.5 to 10 mm in width, with a wave-like () or ruffled pattern and overlapping () arrangement. The tips and edges of these lobes are generally smooth and round, sometimes notched () or incised, often curling up or inward, revealing the paler brown to black underside adorned with hair-like structures () up to 2.5 mm in length. Rhizines are common on the underside of the thallus, except for a brown border near the edges.
The soredia found in this species are and appear white or may become grey due to wear. They are located within specifically structured groups called soralia, which can be linear to oval in shape, often positioned at the edges of the lobes. The presence of soredia causes the lobe margins to curl back and form soralia. The upper surface of the lichen is typically whitish grey to pale greenish-grey, and can be either smooth or slightly wrinkled, without spots (), featuring scattered, shallow cracks. Isidia are absent in this species.
Apothecia (fruiting bodies) are rare in Parmotrema perlatum. When present, they measure 4–8 mm across and are somewhat stalked and funnel-shaped with a brown, concave . The edges of these structures curl inward, becoming thick with soredia as they mature. Its spores are ellipsoid in shape and typically measure between 20 to 28 μm in length and 11 to 17 μm in width, with a wall thickness of 2–3 μm. Pycnidia, which are structures that produce asexual spores called conidia, appear sporadically on the surface () of the thallus, with the conidia being thread-like and straight, measuring 6–8 by 1 μm.
Photobiont
The photobiont partner of Parmotrema perlatum is from Trebouxia, a green algal genus belonging to the order Trebouxiales (order Chlorophyta). It has been identified as an undescribed species within a clade containing Trebouxia arboricola.
A study compared the desiccation tolerance and physiological responses of lichenised Trebouxia to isolated cultures of the same alga. Both forms can survive extended desiccation, but with differing responses to photo-oxidative stress. Lichenisation enhances the photoprotective mechanisms of Trebouxia, improving quenching of excess light energy, particularly under high relative humidity, and controlling reactive oxygen species production under light exposure. However, isolated cultures showed better photosynthetic performance after desiccation recovery. This research demonstrates the mutual benefits of the lichen-photobiont partnership, where the alga gains a sheltered environment boosting its resilience to environmental stressors.
Further studies on Parmotrema perlatum revealed specific antioxidant mechanisms supporting its photobiont under stress. The lichen shows high levels of reactive oxygen species scavenging enzymes such as superoxide dismutase and ascorbate peroxidase, protecting the photobiont from oxidative damage during dehydration and rehydration cycles. This enhanced antioxidant system provides not only physical shelter but also biochemical protection, increasing the photobiont's resilience to environmental fluctuations.
Chemistry
Chemically, Parmotrema perlatum contains atranorin and chloroatranorin, alongside a predominant stictic acid that includes stictic as a major secondary metabolite and smaller amounts of constictic acid and other related substances. Testing the medulla (the inner layer beneath the upper cortex) with spot tests results in K+ (yellow), KC–, and P+ (orange) reactions. The cortical layer, in contrast, is K+ (yellow), KC–, and P–.
The secondary metabolites of Parmotrema perlatum have been studied using gas chromatography–mass spectrometry (GC–MS) and liquid chromatography–mass spectrometry (LC–MS/MS). The lichen produces several notable compounds, including orcinol, atraric acid, benzoic acid, 2,4-dihydroxy-3,6-dimethyl-, methyl ester, and palmitic acid, methyl ester. GC–MS analysis revealed the presence of orcinol (63%) and atraric acid (21%) in the methanol extract, while benzoic acid was predominant in the chloroform extract. The hexane extract contained significant amounts of benzoic acid, 2,4-dihydroxy-3,6-dimethyl-, methyl ester (62%). A more recent study using liquid chromatography-electrospray ionization-mass spectrometry/mass spectrometry as an analysis technique tentatively identified a total of twenty-five lichen products, including 5 depsides, 12 depsidones, 2 diphenyl ethers, 1 aromatic considered as possible artifact, 1 dibenzofuran, 1 carbohydrate, 1 organic acid, and 2 undefined compounds.
Similar species
The distinguishing features of Parmotrema perlatum, such as the presence of soredia and stictic acid, facilitate its easy identification. In mature specimens, the appearance of scattered, fine cracks on the upper surface may resemble the cracked maculae seen in P. reticulatum, which shares similar habitats. However, the two species can be differentiated chemically, as P. reticulatum contains salazinic acid, unlike P. perlatum.
Parmotrema perlatum and Parmotrema stuppeum are two morphologically similar species that can be found in similar habitats. Both species have a loosely attached thallus with revolute, wavy lobes and sparsely ciliate lobe tips. Their upper cortex is continuous and not finely reticulately cracked, while the lower surface is black and rhizinate. Both species also feature linear soralia. However, there are several key differences that can help distinguish between the two. While earlier descriptions suggested that P. stuppeum has a matte, olive-green to brownish-green upper surface and P. perlatum has a slightly shiny, whitish-grey to greyish-green upper surface, recent observations have shown that both species have a distinctly matte upper surface with similar colouration.
The most reliable morphological difference in the field is the location of the soralia: P. stuppeum has strictly terminal soralia, whereas P. perlatum has submarginal soralia. Additionally, the two species can be distinguished by their chemical composition. P. stuppeum contains salazinic acid, while P. perlatum has a stictic acid complex. Although both acids cause a Pd+ orange to orange-red medulla reaction, a potassium (K) spot test can separate the species: the medulla of P. perlatum turns yellow (K+ yellow), whereas in P. stuppeum, the yellow colour turns red (K+ yellow turning red).
Parmotrema perlatum can be distinguished from other sorediate and marginally ciliate species like P. arnoldii and P. robustum by the presence of the stictic acid chemosyndrome. Also, the medulla of P. arnoldii fluoresces strongly when lit with an ultraviolet lamp. Parmotrema perlatum is similar to P. crinitum due to both species having a brown to tan, erhizinate marginal zone and the presence of the stictic acid chemosyndrome in the medulla. However, P. crinitum can be distinguished by its isidiate upper surface. Another potential lookalike, P. margaritatum, is distinguished from P. perlatum by the K+ (red) reaction of its medulla. Cetrelia cetrarioides has been documented as a lookalike, presumably because of the cilia on its thallus margin, and the presence of atranorin and the stictic acid chemosyndrome.
Habitat and distribution
Parmotrema perlatum typically grows in areas with ample light, favouring neutral to slightly acidic-barked broad-leaved trees. It is commonly found on siliceous rocks and walls, as well as mossy coastal rocks, generally growing in places with moderate to strong sunlight. In the Great Smoky Mountains National Park in the United States, Parmotrema perlatum is especially abundant on branches in humid, high-elevation habitats. Similarly, in East Africa, it grows in the misty environments of inselbergs, montane forests, and Erica-dominated habitats, typically found between above sea level.
The species is globally distributed, found in both temperate and tropical regions. It has been reported across numerous European countries including Austria, Belgium, the Czech Republic, France, Germany, Great Britain, Ireland, Italy, Luxembourg, the Netherlands, Portugal, Scandinavia, Slovakia, Spain, and Ukraine. Although it is rare in Eastern Europe, it is widely distributed in both the Asian and European parts of Russia. Beyond Europe, it is also present in Macaronesia, Africa, Australia, North America, and South America. Its Asian distribution includes India, Japan, Taiwan, and South Korea. Although it has historically been recorded in Nepal and Sri Lanka, these reports are considered tentative due to shifting species concepts and possible confusion with the lookalike Parmotrema pseudonilgherrense.
Parmotrema perlatum is globally widespread lichen found on all continents except Antarctica and predominantly in oceanic areas in Europe, primarily grows on bark and occasionally on siliceous rocks amongst mosses. While it is seeing an increase in the Netherlands due to global warming, it is critically endangered in the Czech Republic, Slovakia, and Poland due to susceptibility to air pollution, and is listed as extinct in certain regional Red Data Books due to a lack of recent findings. In contrast, it has been increasing in sightings in the Netherlands, a phenomenon attributed to both global warming decreases in the levels or air pollution in recent decades. Its recent recurrence in Hungary, particularly on some unusual hosts (Catalpa bignonioides, Prunus serotina, and Robinia pseudoacacia) have been suggested as a possible consequence of "a recolonisation process, due to the improving air quality".
Ecology
Parmotrema perlatum is an important species within specific lichen communities in British woodlands, particularly those in late successional mesotrophic settings in oceanic or humid microclimates. It is associated with the Type K Lobaria pulmonaria-Isothecium myosuroides ecological Community. This community type is characterised by its occurrence in mature mesotrophic environments, which are often warmer in winter climates or specific microhabitats. This community includes, in addition to P. perlatum, dominant foliose lichens like Lobaria pulmonaria, Hypotrachyna taylorensis, and Parmotrema crinitum, as well as bryophytes such as Isothecium myosuroides.
A 2017 study investigated the physiological responses of Parmotrema perlatum along an aridity gradient in Southern Portugal. The researchers transplanted thalli of P. perlatum to rural and forested sites characterised by varying levels of aridity and measured several physiological parameters, including photosynthetic performance, pigment content, ergosterol content, and sample viability, both before and after a six-month exposure period. The study found that P. perlatum showed lower photosynthetic performance (measured as FV/FM and the performance index on an absorption basis, PIABS) in drier sites compared to more humid sites. In humid environments, the content of photosynthetic pigments increased post-exposure, while in drier sites, this increase was less pronounced. Additionally, ergosterol content was lower in drier sites, indicating a stress response to arid conditions. These results highlight that P. perlatums physiological responses are significantly influenced by water availability. The ability to maintain higher photosynthetic performance and pigment content in humid conditions suggests that P. perlatum is better adapted to environments with higher moisture levels. This adaptability makes P. perlatum useful as a bioindicator for monitoring ecological responses to climate change and varying moisture conditions in Mediterranean ecosystems.
Lichenicolous (lichen-dwelling) fungi that have been recorded parasitising Parmotrema perlatum include Abrothallus parmotrematis, Briancoppinsia cytospora, Lichenoconium erodens, and Spirographa lichenicola.
Conservation
Parmotrema perlatum has been identified as a species of concern in some regions due to its rarity and declining populations. In Hungary, it has been proposed for 'endangered' status in the Hungarian lichen red list, reflecting its limited distribution and the pressures it faces in its natural habitats there. Similarly, in Ukraine, the species is listed in the Red Data Book of Ukraine with the status of "Rare". Additionally, Parmotrema perlatum is red-listed in Sweden In northern North America, its NatureServe conservation status is designated as "G4", meaning "apparently secure" at the global level. In the United States, it has been assessed as secure in Kentucky and presumed extirpated in Wisconsin, while in Canada, it is considered as vulnerable in British Columbia and Ontario, and critically imperiled in New Brunswick.
Uses
As a spice
Parmotrema perlatum is used as a spice, particularly in the cuisine of Tamil Nadu. It is especially prevalent in Chettinad cuisine, being used in the popular rice dish biryani, and also in many meat and vegetarian dishes.
In its raw state, black stone flower does not have much taste or fragrance. However, when put in contact with heat, especially hot cooking oil or ghee, it releases a distinctive earthy, smoky flavour and aroma. This property of black stone flower is especially valued in the tempering step of cooking a number of Indian dishes.
The spice is also integral to various regional masalas throughout the Indian subcontinent. Parmotrema perlatum is a key ingredient in masalas such as Kala and Goda masala of Maharashtra, Anglo-Indian bottle masala, bhojwar masala from Hyderabad, and potli masala in Lucknow. It is what many cooks and commercial spice blend makers believe sets apart accomplished dishes from those made by amateurs. Despite its lack of a specific aroma or describable flavour in its raw form, its contribution to the complex flavour profile of these spice blends is highly valued.
Dyeing
A natural purple dye extracted from Parmotrema perlatum using ammonia fermentation showed optimal results, with a notable dye yield and effective application on silk fabric. The study demonstrated the dye's potential as a sustainable alternative to synthetic dyes, with satisfactory colour fastness and fabric strength enhancement.
Recent research highlights the antimicrobial, antioxidant, and photocatalytic capabilities of zinc oxide nanoparticles synthesised using Parmotrema perlatum, marking a significant step towards sustainable dyeing practices and broadening the lichen's applicative horizons.
Traditional medicines
Parmotrema perlatum is used as a component of a herbal mixture in Ayurvedic medicine, one of several lichen species used as charila. Referenced in ancient Ayurvedic texts and first mentioned in the Atharvaveda around 1500 BCE, charila is a lichen mixture traditionally used in India for its purported medicinal properties. It has been employed to treat various ailments, including digestive and respiratory issues, skin conditions, and reproductive health concerns, and it also serves as an ingredient in treatments for infertility. For chronic ulcers, a powder made from dried lichen, infused in pork suet, is applied externally.
Biomonitoring
Parmotrema perlatum is sensitive to air pollution, making it a useful bioindicator. This sensitivity is utilised in the "Hawksworth and Rose" scale, which estimates mean winter sulphur dioxide (SO2) levels in England and Wales by observing lichens on acidic and nutrient-poor bark. According to this scale, P. perlatum is found only in zones 8 to 10, indicating areas with the lowest SO2 concentrations, less than 35 micrograms per cubic metre (μg/m3).
A 2022 study analyzed the effects of SO2 and nitrogen dioxide (NO2) fumigation on the chlorophyll content of Parmotrema perlatum collected from the Mount Lawu volcano in Indonesia. The results indicated that increased exposure to these pollutants leads to a significant reduction in chlorophyll levels. The study demonstrated that SO2 and NO2 negatively impact the physiological processes of the lichen, particularly its photosynthetic efficiency, demonstrating the sensitivity of P. perlatum to air pollution.
Parmotrema perlatum has been effectively used in biomonitoring studies to assess environmental radioactivity. Research conducted in Turkey found that this lichen species retains radioactive caesium-137 (137Cs) from atmospheric deposition, such as fallout from the Chernobyl accident. The ecological half-life of 137Cs in Parmotrema perlatum was determined to be approximately 5.5 years, indicating its capability to monitor long-term radioactive contamination in the environment.
Research
Research on the bioactive properties of Parmotrema perlatum has revealed several findings. The methanol extract of this species has been shown to significantly reduce blood glucose levels in streptozotocin-induced diabetic rats, attributed to its inhibitory activity on alpha-glucosidase rather than an effect on insulin secretion. This extract also has a high phenolic content and moderate antioxidant capacity, which could help prevent secondary complications of diabetes. The antioxidant potential and free radical-scavening activity of P. perlatum extracts has been further demonstrated through various chemical assays.
Additionally, Parmotrema perlatum has some antimicrobial properties. The crude polysaccharide fraction of this lichen demonstrated antibacterial activity against Escherichia coli and Staphylococcus aureus, which are common pathogens in diabetic foot ulcers. Furthermore, extracts from this species showed significant antiviral activity against the yellow fever virus envelope. Tests against the Gram-negative bacteria Pseudomonas aeruginosa, Chromobacterium violaceum, and Gram-positive Lactobacillus plantarum showed that the methanol extract had the highest antibacterial activity among the three solvent extracts evaluated.
In terms of cytotoxic and anticancer activities, the n-hexane, diethyl ether, and methanol extracts of Parmotrema perlatum have been studied against various cancer cell lines, with the n-hexane extract showing the highest cytotoxic effects. The extracts were particularly effective against murine Lewis lung carcinoma and human glioblastoma cell lines.
| Biology and health sciences | Lichens | Plants |
27046146 | https://en.wikipedia.org/wiki/Arithmetic%20logic%20unit | Arithmetic logic unit | In computing, an arithmetic logic unit (ALU) is a combinational digital circuit that performs arithmetic and bitwise operations on integer binary numbers. This is in contrast to a floating-point unit (FPU), which operates on floating point numbers. It is a fundamental building block of many types of computing circuits, including the central processing unit (CPU) of computers, FPUs, and graphics processing units (GPUs).
The inputs to an ALU are the data to be operated on, called operands, and a code indicating the operation to be performed; the ALU's output is the result of the performed operation. In many designs, the ALU also has status inputs or outputs, or both, which convey information about a previous operation or the current operation, respectively, between the ALU and external status registers.
Signals
An ALU has a variety of input and output nets, which are the electrical conductors used to convey digital signals between the ALU and external circuitry. When an ALU is operating, external circuits apply signals to the ALU inputs and, in response, the ALU produces and conveys signals to external circuitry via its outputs.
Data
A basic ALU has three parallel data buses consisting of two input operands (A and B) and a result output (Y). Each data bus is a group of signals that conveys one binary integer number. Typically, the A, B and Y bus widths (the number of signals comprising each bus) are identical and match the native word size of the external circuitry (e.g., the encapsulating CPU or other processor).
Opcode
The opcode input is a parallel bus that conveys to the ALU an operation selection code, which is an enumerated value that specifies the desired arithmetic or logic operation to be performed by the ALU. The opcode size (its bus width) determines the maximum number of distinct operations the ALU can perform; for example, a four-bit opcode can specify up to sixteen different ALU operations. Generally, an ALU opcode is not the same as a machine language instruction, though in some cases it may be directly encoded as a bit field within such instructions.
Status
Outputs
The status outputs are various individual signals that convey supplemental information about the result of the current ALU operation. General-purpose ALUs commonly have status signals such as:
Carry-out, which conveys the carry resulting from an addition operation, the borrow resulting from a subtraction operation, or the overflow bit resulting from a binary shift operation.
Zero, which indicates all bits of Y are logic zero.
Negative, which indicates the result of an arithmetic operation is negative.
Overflow, which indicates the result of an arithmetic operation has exceeded the numeric range of Y.
Parity, which indicates whether an even or odd number of bits in Y are logic one.
Inputs
The status inputs allow additional information to be made available to the ALU when performing an operation. Typically, this is a single "carry-in" bit that is the stored carry-out from a previous ALU operation.
Circuit operation
An ALU is a combinational logic circuit, meaning that its outputs will change asynchronously in response to input changes. In normal operation, stable signals are applied to all of the ALU inputs and, when enough time (known as the "propagation delay") has passed for the signals to propagate through the ALU circuitry, the result of the ALU operation appears at the ALU outputs. The external circuitry connected to the ALU is responsible for ensuring the stability of ALU input signals throughout the operation, and for allowing sufficient time for the signals to propagate through the ALU circuitry before sampling the ALU outputs.
In general, external circuitry controls an ALU by applying signals to the ALU inputs. Typically, the external circuitry employs sequential logic to generate the signals that control ALU operation. The external sequential logic is paced by a clock signal of sufficiently low frequency to ensure enough time for the ALU outputs to settle under worst-case conditions (i.e., conditions resulting in the maximum possible propagation delay).
For example, a CPU starts an addition operation by routing the operands from their sources (typically processor registers) to the ALU's operand inputs, while simultaneously applying a value to the ALU's opcode input that configures it to perform an addition operation. At the same time, the CPU enables the destination register to store the ALU output (the resulting sum from the addition operation) upon operation completion. The ALU's input signals, which are held stable until the next clock, are allowed to propagate through the ALU and to the destination register while the CPU waits for the next clock. When the next clock arrives, the destination register stores the ALU result and, since the ALU operation has completed, the ALU inputs may be set up for the next ALU operation.
Functions
A number of basic arithmetic and bitwise logic functions are commonly supported by ALUs. Basic, general purpose ALUs typically include these operations in their repertoires:
Arithmetic operations
Add: A and B are summed and the sum appears at Y and carry-out.
Add with carry: A, B and carry-in are summed and the sum appears at Y and carry-out.
Subtract: B is subtracted from A (or vice versa) and the difference appears at Y and carry-out. For this function, carry-out is effectively a "borrow" indicator. This operation may also be used to compare the magnitudes of A and B; in such cases the Y output may be ignored by the processor, which is only interested in the status bits (particularly zero and negative) that result from the operation.
Subtract with borrow: B is subtracted from A (or vice versa) with borrow (carry-in) and the difference appears at Y and carry-out (borrow out).
Two's complement: A (or B) is subtracted from zero and the difference appears at Y.
Increment: A (or B) is increased by one and the resulting value appears at Y.
Decrement: A (or B) is decreased by one and the resulting value appears at Y.
Bitwise logical operations
AND: the bitwise AND of A and B appears at Y.
OR: the bitwise OR of A and B appears at Y.
Exclusive-OR: the bitwise XOR of A and B appears at Y.
Ones' complement: all bits of A (or B) are inverted and appear at Y.
Bit shift operations
ALU shift operations cause operand A (or B) to shift left or right (depending on the opcode) and the shifted operand appears at Y. Simple ALUs typically can shift the operand by only one bit position, whereas more complex ALUs employ barrel shifters that allow them to shift the operand by an arbitrary number of bits in one operation. In all single-bit shift operations, the bit shifted out of the operand appears on carry-out; the value of the bit shifted into the operand depends on the type of shift.
Arithmetic shift: the operand is treated as a two's complement integer, meaning that the most significant bit is a "sign" bit and is preserved.
Logical shift: a logic zero is shifted into the operand. This is used to shift unsigned integers.
Rotate: the operand is treated as a circular buffer of bits in which its least and most significant bits are effectively adjacent.
Rotate through carry: the carry bit and operand are collectively treated as a circular buffer of bits.
Other operations
Pass through: all bits of A (or B) appear unmodified at Y. This operation is typically used to determine the parity of the operand or whether it is zero or negative, or to copy the operand to a processor register.
Applications
Status usage
Upon completion of each ALU operation, the ALU's status output signals are usually stored in external registers to make them available for future ALU operations (e.g., to implement multiple-precision arithmetic) and for controlling conditional branching. The bit registers that store the status output signals are often collectively treated as a single, multi-bit register, which is referred to as the "status register" or "condition code register".
Depending on the ALU operation being performed, some status register bits may be changed and others may be left unmodified. For example, in bitwise logical operations such as AND and OR, the carry status bit is typically not modified as it is not relevant to such operations.
In CPUs, the stored carry-out signal is usually connected to the ALU's carry-in net. This facilitates efficient propagation of carries (which may represent addition carries, subtraction borrows, or shift overflows) when performing multiple-precision operations, as it eliminates the need for software-management of carry propagation (via conditional branching, based on the carry status bit).
Multiple-precision arithmetic
In integer arithmetic computations, multiple-precision arithmetic is an algorithm that operates on integers which are larger than the ALU word size. To do this, the algorithm treats each integer as an ordered collection of ALU-size fragments, arranged from most-significant (MS) to least-significant (LS) or vice versa. For example, in the case of an 8-bit ALU, the 24-bit integer 0x123456 would be treated as a collection of three 8-bit fragments: 0x12 (MS), 0x34, and 0x56 (LS). Since the size of a fragment exactly matches the ALU word size, the ALU can directly operate on this "piece" of operand.
The algorithm uses the ALU to directly operate on particular operand fragments and thus generate a corresponding fragment (a "partial") of the multi-precision result. Each partial, when generated, is written to an associated region of storage that has been designated for the multiple-precision result. This process is repeated for all operand fragments so as to generate a complete collection of partials, which is the result of the multiple-precision operation.
In arithmetic operations (e.g., addition, subtraction), the algorithm starts by invoking an ALU operation on the operands' LS fragments, thereby producing both a LS partial and a carry out bit. The algorithm writes the partial to designated storage, whereas the processor's state machine typically stores the carry out bit to an ALU status register. The algorithm then advances to the next fragment of each operand's collection and invokes an ALU operation on these fragments along with the stored carry bit from the previous ALU operation, thus producing another (more significant) partial and a carry out bit. As before, the carry bit is stored to the status register and the partial is written to designated storage. This process repeats until all operand fragments have been processed, resulting in a complete collection of partials in storage, which comprise the multi-precision arithmetic result.
In multiple-precision shift operations, the order of operand fragment processing depends on the shift direction. In left-shift operations, fragments are processed LS first because the LS bit of each partial—which is conveyed via the stored carry bit—must be obtained from the MS bit of the previously left-shifted, less-significant operand. Conversely, operands are processed MS first in right-shift operations because the MS bit of each partial must be obtained from the LS bit of the previously right-shifted, more-significant operand.
In bitwise logical operations (e.g., logical AND, logical OR), the operand fragments may be processed in any arbitrary order because each partial depends only on the corresponding operand fragments (the stored carry bit from the previous ALU operation is ignored).
Complex operations
Although it is possible to design ALUs that can perform complex functions, this is usually impractical due to the resulting increases in circuit complexity, power consumption, propagation delay, cost and size. Consequently, ALUs are typically limited to simple functions that can be executed at very high speeds (i.e., very short propagation delays), with more complex functions being the responsibility of external circuitry. For example:
In simple cases in which a CPU contains a single ALU, the CPU typically implements a complex operation by orchestrating a sequence of ALU operations according to a software algorithm.
More specialized architectures may use multiple ALUs to accelerate complex operations. In such systems, the ALUs are often pipelined, with intermediate results passing through ALUs arranged like a factory production line. Performance is greatly improved over that of a single ALU because all of the ALUs operate concurrently and software overhead is significantly reduced.
Graphics processing units
Graphics processing units (GPUs) often contain hundreds or thousands of ALUs which can operate concurrently. Depending on the application and GPU architecture, the ALUs may be used to simultaneously process unrelated data or to operate in parallel on related data. An example of the latter is graphics rendering, in which multiple ALUs perform the same operation in parallel on a group of pixels, with each ALU operating on a pixel within a scene.
Intel use ALU term to equal to the GPU shader.
Implementation
An ALU is usually implemented either as a stand-alone integrated circuit (IC), such as the 74181, or as part of a more complex IC. In the latter case, an ALU is typically instantiated by synthesizing it from a description written in VHDL, Verilog or some other hardware description language. For example, the following VHDL code describes a very simple 8-bit ALU:
entity alu is
port ( -- the alu connections to external circuitry:
A : in signed(7 downto 0); -- operand A
B : in signed(7 downto 0); -- operand B
OP : in unsigned(2 downto 0); -- opcode
Y : out signed(7 downto 0)); -- operation result
end alu;
architecture behavioral of alu is
begin
case OP is -- decode the opcode and perform the operation:
when "000" => Y <= A + B; -- add
when "001" => Y <= A - B; -- subtract
when "010" => Y <= A - 1; -- decrement
when "011" => Y <= A + 1; -- increment
when "100" => Y <= not A; -- 1's complement
when "101" => Y <= A and B; -- bitwise AND
when "110" => Y <= A or B; -- bitwise OR
when "111" => Y <= A xor B; -- bitwise XOR
when others => Y <= (others => 'X');
end case;
end behavioral;
History
Mathematician John von Neumann proposed the ALU concept in 1945 in a report on the foundations for a new computer called the EDVAC.
The cost, size, and power consumption of electronic circuitry was relatively high throughout the infancy of the Information Age. Consequently, all early computers had a serial ALU that operated on one data bit at a time although they often presented a wider word size to programmers. The first computer to have multiple parallel discrete single-bit ALU circuits was the 1951 Whirlwind I, which employed sixteen such "math units" to enable it to operate on 16-bit words.
In 1967, Fairchild introduced the first ALU-like device implemented as an integrated circuit, the Fairchild 3800, consisting of an eight-bit arithmetic unit with accumulator. It only supported adds and subtracts but no logic functions.
Full integrated-circuit ALUs soon emerged, including four-bit ALUs such as the Am2901 and 74181. These devices were typically "bit slice" capable, meaning they had "carry look ahead" signals that facilitated the use of multiple interconnected ALU chips to create an ALU with a wider word size. These devices quickly became popular and were widely used in bit-slice minicomputers.
Microprocessors began to appear in the early 1970s. Even though transistors had become smaller, there was sometimes insufficient die space for a full-word-width ALU and, as a result, some early microprocessors employed a narrow ALU that required multiple cycles per machine language instruction. Examples of this includes the popular Zilog Z80, which performed eight-bit additions with a four-bit ALU. Over time, transistor geometries shrank further, following Moore's law, and it became feasible to build wider ALUs on microprocessors.
Modern integrated circuit (IC) transistors are orders of magnitude smaller than those of the early microprocessors, making it possible to fit highly complex ALUs on ICs. Today, many modern ALUs have wide word widths, and architectural enhancements such as barrel shifters and binary multipliers that allow them to perform, in a single clock cycle, operations that would have required multiple operations on earlier ALUs.
ALUs can be realized as mechanical, electro-mechanical or electronic circuits and, in recent years, research into biological ALUs has been carried out (e.g., actin-based).
| Technology | Computer hardware | null |
2728240 | https://en.wikipedia.org/wiki/Carnallite | Carnallite | Carnallite (also carnalite) is an evaporite mineral, a hydrated potassium magnesium chloride with formula KCl.MgCl2·6(H2O). It is variably colored yellow to white, reddish, and sometimes colorless or blue. It is usually massive to fibrous with rare pseudohexagonal orthorhombic crystals. The mineral is deliquescent (absorbs moisture from the surrounding air) and specimens must be stored in an airtight container.
Carnallite occurs with a sequence of potassium and magnesium evaporite minerals: sylvite, kainite, picromerite, polyhalite, and kieserite. Carnallite is an uncommon double chloride mineral that only forms under specific environmental conditions in an evaporating sea or sedimentary basin. It is mined for both potassium and magnesium and occurs in the evaporite deposits of Carlsbad, New Mexico; the Paradox Basin in Colorado and Utah; Stassfurt, Germany; the Perm Basin, Russia; and the Williston Basin in Saskatchewan, Canada. These deposits date from the Devonian through the Permian Periods. In contrast, both Israel and Jordan produce potash from the Dead Sea by using evaporation pans to further concentrate the brine until carnallite precipitates, dredging the carnallite from the pans, and processing to remove the magnesium chloride from the potassium chloride.
Carnallite was first described in 1856 from its type location of Stassfurt Deposit, Saxony-Anhalt, Germany. It was named for the Prussian mining engineer Rudolf von Carnall (1804–1874).
Background
Halides are binary compounds. They are composed of a halogen and a metal ion. The crystal chemistry of halides is characterized by the electronegativity of halogen ions. This means that the dominant large ions are the Cl−, Br−, F−, or I−. These are easily polarized. The ions combine with similarly large but low valence and weakly polarized cations. The cations are mostly of the alkali metal group. Sylvite is a binary compound with the formula KCl. Sylvite precipitates first from mixed solutions of K+, Mg2+ and Cl−, leaving a brine enriched in magnesium from which the mixed halide carnallite then precipitates.
Composition
Carnallite's chemical formula is KMgCl3·6(H2O). Synthetic carnallite crystal specimens can be produced from 1.5 mole percent KCl and 98.5 mole percent MgCl2·6H2O by slow crystallization at 25 °C. Its density is 1.602 g/cm3. Carnallite can also be produced by grinding the combination of hydrated magnesium chloride and potassium chloride.
Structure
The carnallite structure exhibits corner- and face-sharing. There is a network of KCl6 octahedra, with two-thirds of them sharing faces. Mg(H2O)6 octahedra occupy the open spaces within the KCl octahedra. The interatomic distance between Mg and H2O ranges from 0.204 to 0.209 nm, with an average is 0.2045 nm. The interatomic distance between K and Cl ranges 0.317 to 0.331 nm., with an average of 0.324 nm. The resulting structure has a calculated density of 1.587 g/cm3, in good agreement with the measured value of 1.602 g/cm3.
Face-sharing creates more chance of instability, according to the third of Pauling's rules. In carnallite, the water molecules enclose the magnesium ions. This prevents the magnesium and the chloride from interacting directly; instead, the water molecules act as charge transmitters. The five chloride anions are each coordinated to two potassium cations as well as four water molecules. This means that each chloride anion receives 1/6 of a +1 charge from each of the two potassium ions. The chloride also obtains 1/6 of a +1 charge from each the four water molecules. The charges thus total six 1/6 positive charges, which balance the negative charge of the chloride. These two aspects render the rare face-sharing described by the second and third of Pauling's rules acceptable in the carnallite structure.
Physical properties
Carnallite's refractive index ranges from 1.467 to 1.494. Carnallite may be red as a result of hematite (Fe2O3) inclusions. The fragmented shards of iron oxide produce red tints in the thin laminae of hematite. Carnallite is also deliquescent in high humidity. This implies that it is also extremely soluble in water. Individual crystals are pseudo-hexagonal and tabular but are extremely rarely seen. Field indicators of carnallite are environment of formation, absence of cleavage, and fracture. Other indicators can be density, taste, associations to local minerals, and whether it is capable of luminescence. Carnallite has a bitter taste. Carnallite may not only be fluorescent but is capable of being phosphorescent. The potassium that carnallite contains fuses easily within a flame, creating a violet color.
Geologic occurrence
Mineral associations based on some physical properties include, but not limited to, halite, anhydrite, dolomite, gypsum, kainite, kieserite, polyhalite, and sylvite.
Carnallite minerals are mineral sediments known as evaporites. Evaporites are concentrated by evaporation of seawater. The inflow of water must be below the evaporation or use levels. This creates a prolonged evaporation period. In controlled environment experiments, the halides form when 10%–20% of the original sample of water remains. Closer to 10 percent sylvite followed by Carnallite form.
Carnallite is mostly found in saline marine deposits, although beds exist in the endorheic Qaidam Basin of China's Qinghai Province near Dabusun Nor.
Uses
Carnallite is mostly used in fertilizers. It is an important source of potash. Only sylvite outranks carnallite's importance in potash production. Both are uncommon because they are some of the last evaporites to form. Soluble potassium salts are the main sources for fertilizer. This is because the potassium is difficult to separate from insoluble potassium feldspar. Carnallite is a minor source of magnesium worldwide; however, it is Russia's main source.
| Physical sciences | Sedimentary rocks | Earth science |
2729585 | https://en.wikipedia.org/wiki/Working%20animal | Working animal | A working animal is an animal, usually domesticated, that is kept by humans and trained to perform tasks instead of being slaughtered to harvest animal products. Some are used for their physical strength (e.g. oxen and draft horses) or for transportation (e.g. riding horses and camels), while others are service animals trained to execute certain specialized tasks (e.g. hunting and guide dogs, messenger pigeons, and fishing cormorants). They may also be used for milking or herding. Some, at the end of their working lives, may also be used for meat or leather.
The history of working animals may predate agriculture as dogs were used by hunter-gatherer ancestors; around the world, millions of animals work in relationship with their owners. Domesticated species are often bred for different uses and conditions, especially horses and working dogs. Working animals are usually raised on farms, though some are still captured from the wild, such as dolphins and some Asian elephants.
People have found uses for a wide variety of abilities in animals, and even industrialized societies use many animals for work. People use the strength of horses, elephants, and oxen to pull carts and move loads. Police forces use dogs for finding illegal substances and assisting in apprehending wanted persons, others use dogs to find game or search for missing or trapped people. People use various animals—camels, donkeys, horses, dogs, etc.—for transport, either for riding or to pull wagons and sleds. Other animals, including dogs and monkeys, help disabled people.
On rare occasions, wild animals are not only tamed, but trained to perform work—though often solely for novelty or entertainment, as such animals tend to lack the trustworthiness and mild temper of true domesticated working animals. Conversely, not all domesticated animals are working animals. For example, while cats may catch mice, it is an instinctive behavior, not one that can be trained by human intervention. Other domesticated animals, such as sheep or rabbits, may have agricultural uses for meat, hides and wool, but are not suitable for work. Finally, small domestic pets, such as most small birds (other than certain types of pigeon) are generally incapable of performing work other than providing companionship.
Roles and specializations
Transportation
Some animals are used due to sheer physical strength in tasks such as ploughing or logging. Such animals are grouped as a draught or draft animals. Others may be used as pack animals, for animal-powered transport, the movement of people and goods. Together, these are sometimes called beasts of burden. Some animals are ridden by people on their backs and are known as mounts. Alternatively, one or more animals in harness may be used to pull vehicles.
Riding animals or mounts
Riding animals are animals that people use as mounts in order to perform tasks such as traversing across long distances or over rugged terrain, hunting on horseback or with some other riding animal, patrolling around rural and/or wilderness areas, rounding up and/or herding livestock or even for recreational enjoyment. They mainly include equines such as horses, donkeys, and mules; bovines such as cattle, water buffalo, and yak. In some places, elephants, llamas and camels are also used. Dromedary camels are in arid areas of Australia, North Africa and the Middle East; the less common Bactrian camel inhabits central and East Asia; both are used as working animals. On occasion, reindeer, though usually driven, may be ridden.
Certain wild animals have been tamed and used for riding, usually for novelty purposes, including the zebra and the ostrich. Some mythical creatures are believed to act as divine mounts, such as garuda in Hinduism (See vahana for divine mounts in Hinduism) and the winged horse Pegasus in Greek mythology.
Pack animals
Pack animals may be of the same species as mounts or harness animals, though animals such as horses, mules, donkeys, reindeer and both types of camel may have individual bloodlines or breeds that have been selectively bred for packing. Additional species are only used to carry loads, including llamas in the Andes.
Domesticated cattle and yaks are also used as pack animals. Other species used to carry cargo include dogs and pack goats.
Draft animals
An intermediate use is as draft animals, harnessed singly or in teams, to pull sleds, wheeled vehicles or ploughs.
Oxen are slow but strong, and have been used in a yoke since ancient times: the earliest surviving vehicle, Puabi's Sumerian sledge, was ox-drawn; an acre was originally defined as the area a span of oxen could plow in a day. The domestic water buffalo and carabao, pull wagons and ploughs in Southeast Asia and the Philippines.
Draught or draft horses are commonly used in harness for heavy work. Several breeds of medium-weight horses are used to pull lighter wheeled carts, carriages and buggies when a certain amount of speed or style is desirable.
Mules are considered tough and strong, with harness capacity dependent on the type of horse mare used to produce the mule foal. Because they are a hybrid animal and usually are infertile, separate breeding programs must also be maintained.
Ponies and donkeys are often used to pull carts and small wagons. Historically, ponies were commonly used in mining to pull ore carts.
Dogs are used for pulling light carts or, particularly, sleds (e.g. sled dogs such as huskies) for both recreation and working purposes.
Goats also can perform light harness work in front of carts.
Reindeer are used in the Arctic and sub-Arctic Nordic countries and Siberia. During World War II, the Red Army deployed deer transportation battalions on the Eastern Front. In the twenty-first century, Russian soldiers continue to train with reindeer sleds in winter. In traditional festive legend, Santa Claus's reindeer pull a sleigh through the night sky to help Santa Claus deliver gifts to children on Christmas Eve.
Elephants are used for logging in Southeast Asia.
Less often, camels and llamas have been trained to harness. According to Juan Ignacio Molina the Dutch captain Joris van Spilbergen observed the use of chiliquenes (a llama type) by native Mapuches of Mocha Island as plough animals in 1614.
Assorted wild animals have, on occasion, been tamed and trained to harness, including zebras and even moose.
Guard animals
As some domesticated animals display extremely protective or territorial behavior, certain breeds and species have been utilized to guard people and/or property such as homes, public buildings, businesses, crops, livestock and even venues of criminal activity. Guard animals can either act as alarms to alert their owners of danger or they can be used to actively scare off and/or even attack encroaching intruders or dangerous animals. Well known examples of guard animals include dogs, geese and llamas.
Powering fixed machinery
Working draught animals may power fixed machinery using a treadmill and have been used throughout history to power a winch to raise water from a well. Turnspit dogs were formerly used to power roasting jacks for roasting meat.
Treatment animals
Working as a form of biological treatment for the environment. Animals such as Asian carps were imported to the U.S. in 1970s to control algae, weed, and parasite growth in aquatic farms, weeds in canal systems, and as one form of sewage treatment.
Pathogens and diseases
Animals can be used to detect the presence of pathogens and patients carrying infectious diseases.
Dogs (including scent hounds) and bees have been trained to detect COVID-19 infections.
Dogs have been trained to detect cancer. One study reported ants could be used to detect cancer via urine.
Detection rats such as those trained by APOPO can also be taught to identify diseases, especially pulmonary tuberculosis.
Searching and retrieving
Dogs and pigs, with a better sense of smell than humans, can assist with gathering by finding valuable products, such as truffles (a very expensive subterranean fungus). The French typically use truffle hogs, while Italians mainly use dogs.
Monkeys are trained to pick coconuts from palm trees, a job many human workers consider as too dangerous.
Detecting contraband
Detection dogs, commonly employed by law enforcement authorities, are trained to use their senses to detect illegal drugs, explosives, currency, and contraband electronics such as illicit mobile phones, among other things. The sense most used by detection dogs is smell, hence such dogs are also commonly known as 'sniffer dogs'. For this task, dogs may sometimes be used remotely from the suspect item, for example via the Remote Air Sampling for Canine Olfaction (RASCO) system.
Interfacing and organization
Assistance animals
The best-known example is the guide dog or seeing eye dog for blind people. | Technology | Agriculture, labor and economy | null |
5006326 | https://en.wikipedia.org/wiki/Salix%20nigra | Salix nigra | Salix nigra, the black willow, is a species of willow native to a large portion of North America, from New Brunswick and southern Ontario west to Arizona and California, and south to northern Florida and Texas.
Description
Salix nigra is a medium-sized deciduous tree, the largest North American species of willow, growing to tall, exceptionally up to , with a trunk diameter. The bark is dark brown to blackish, becoming fissured in older trees, and frequently forking near the base. The shoots are slender and variable in color from green to brown, yellow or purplish; they are (like the related European Salix fragilis) brittle at the base, snapping evenly at the branch junction if bent sharply. The foliage buds are long, with a single, pointed reddish-brown bud scale. The leaves are alternate, long, thin, long and broad, usually somewhat falcate, dark, shiny green on both sides or with a lighter green underside, with a finely serrated margin, a short petiole and a pair of small stipules. It is dioecious, with small, greenish yellow to yellow flowers borne on catkins long in early spring at the same time as the new leaves appear. The fruit is a capsule which splits open when mature to release the numerous minute, down-covered seeds. The leaves turn a lemon yellow in the fall.
Distribution and habitat
Salix nigra is native to a large portion of North America, from New Brunswick and southern Ontario west to California and Arizona, and south to northern Florida and Texas. It is also found in parts of Mexico, both south and west of the Rio Grande. It has also been introduced along streams in the state of Utah. Salix nigra grows best in areas of full sun and wet or moist soils. Thus, it is typically found along streams and in swamps.
Taxonomy
Black willow is part of the Salicaceae, the willow family. The accepted name for black willow is: Salix nigra Marshall. Marshall, the "Father of American Dendrology", first described this taxon in 1785.
Salix gooddingii (Goodding's willow) is sometimes considered a variety of S. nigra as S. nigra var. vallicola Dudley; when recognized, this extends the range of S. nigra to western North America. However, the two are usually treated as distinct species. Some other related taxa and synonyms are S. nigra var. altissima, S. nigra var. brevijulis, S. nigra var. longifolia, S. nigra var. marginata, and S. nigra var. wardii.
Another name occasionally used for black willow is "swamp willow", not to be confused with Salix myrtilloides (swamp willow). Other common names include "Goodding willow", "southwestern black willow", "Dudley willow", and "sauz" (a Spanish word).
Reproduction and growth
Salix nigra is dioecious, which means it has separate male and female trees. Flowering may be climate dependent. It flowers during February in the southern part of its range, and flowers until June in the northern parts. These trees are capable of producing seeds when they are around the age of 10 years. The black willow continuously has good seed crops year after year, with only a few failures. The seeds require very specific conditions to germinate. They prefer soil that is very wet or flooded. After they germinate, they can experience excellent growth if they are exposed to high sunlight and copious moisture during the growing season.
Black willow are capable of developing special features related to flood tolerance to help them survive in flooded habitats. Under flooded conditions, black willow develop hypertrophied lenticels and water roots. Black willows are also sensitive to drought conditions. Black willows living in drought conditions experience inhibition of their branch and root growth.
Early season leaves of some species in the Salix genus generally contain denser indumentum than leaves that are produced later in the season. Indumentum refers to hairs that are red or brown in color. Salix have first leaves, which are leaves that grow from the bud, and can also be called leaves that are "preformed". The leaves that form as the branch continues to grow out are called new leaves, or "neoformed".
Largest example
According to the National Register of Big Trees, the largest black willow tree in the US is in Hennepin, Minnesota. Its height is , circumference is and spread is .
The Marlboro Tree, located in Marlboro Township, New Jersey is certified by the State of New Jersey as the largest known example of this tree in the state. It is about 152 years old and measures in height and in circumference. Five grown people must hold hands to fully encircle the tree.
Uses
Black willow roots are very bitter, and have been used as a substitute for quinine in the past. Ethnobotanical uses of black willow by various Native American tribes include basketry, and treatment of fever, headache, and coughs. It was recognized that using the bark and leaves of Salix nigra was useful in treating rheumatism.
The black willow is the only United States native willow species to be used as timber for a variety of different items. Black willow lumber is used in furniture and shipping containers. The largest production site for black willow timber was in Louisiana at its peak during the 1970s.
The wood of Salix nigra is very lightweight. The wood was once used for artificial limbs, such as wooden hands. It is also capable of maintaining its shape, does not splinter very easily, and has a moderately high shock resistance, allowing it to sustain continuous moderate impacts.
It may also be used in environmental restoration. Black willow is very resistant to herbivory, flooding, and is an erosion control tool. Salix nigra is used for marshland stabilization or restoration projects as long as the roots don't penetrate any clay liner that may be in place. Salix nigra demonstrated some success as a photodegradation tool. Photodegradation is theorized to work by drawing the target chemical, such as bentazon, up the roots and stem into the leaves where it is degraded by high-energy radiation provided by the sun.
Disease
Black willows are susceptible to diseases such as crown gall and cankers. Crown gall is caused by bacteria living in the soil where the black willow is present. When black willow are infected, the bacteria stimulate a quick burst in growth of plant cells. They cause the tree to form tumor-like growths, or "galls" on different parts of the tree such as their roots or on the lower branches. As these galls get bigger, they become hard, woody, brown in color, and corky. Black willows do not normally die from crown galls. However, these galls can cause a disruption in the flow of nutrients throughout the tree and can have its normal growth stunted or slowed down due to this disease. Crown galls are less problematic in soils that are more acidic; thus, soil pH can be an important factor in helping to limit this disease.
Black canker disease is caused by the fungus Glomerella miyabeana. This disease starts as spotting on the leaves of the black willow, where it then spreads to the petiole and eventually the twig or branch holding the leaf. It then forms black patches on the stem of the tree that are capable of expanding. Black willows that are in nutrient-poor locations or that are experiencing poor temporary climate conditions (such as a temporary lack of rainfall or short, drought-like conditions) are considered stressed and very susceptible to the spread of these cankers. Tissue on the black willow that is affected by these cankers will not grow with the rest of the tree, and cracks will begin to form.
| Biology and health sciences | Malpighiales | Plants |
5009517 | https://en.wikipedia.org/wiki/Sea%20snake | Sea snake | Sea snakes, or coral reef snakes, are elapid snakes that inhabit marine environments for most or all of their lives. They belong to two subfamilies, Hydrophiinae and Laticaudinae. Hydrophiinae also includes Australasian terrestrial snakes, whereas Laticaudinae only includes the sea kraits (Laticauda), of which three species are found exclusively in freshwater. If these three freshwater species are excluded, there are 69 species of sea snakes divided among seven genera.
Most sea snakes are venomous, except the genus Emydocephalus, which feeds almost exclusively on fish eggs. Sea snakes are extensively adapted to a fully aquatic life and are unable to move on land, except for the sea kraits, which have limited land movement. They are found in warm coastal waters from the Indian Ocean to the Pacific and are closely related to venomous terrestrial snakes in Australia.
All sea snakes have paddle-like tails and many have laterally compressed bodies that give them an eel-like appearance. Unlike fish, they do not have gills and must surface regularly to breathe. Along with Cetaceans, they are among the most completely aquatic of all extant air-breathing vertebrates. Among this group are species with some of the most potent venoms of all snakes. Some have gentle dispositions and bite only when provoked, while others are much more aggressive.
Description
The majority of adult sea snakes species grow to between in length, with the largest, Hydrophis spiralis, reaching a maximum of . Their eyes are relatively small with a round pupil and most have nostrils located dorsally. The skulls do not differ significantly from those of terrestrial elapids, although their dentition is relatively primitive with short fangs and (with the exception of Emydocephalus) as many as 18 smaller teeth behind them on the maxilla.
Most sea snakes are completely aquatic and have adapted to sea environments in many ways, the most characteristic of which is a paddle-like tail that has improved their swimming ability. To a varying degree, the bodies of many species are laterally compressed, especially in the pelagic species. This has often caused the ventral scales to become reduced in size, even difficult to distinguish from the adjoining scales. Their lack of ventral scales means they have become virtually helpless on land, but as they live out their entire lifecycles at sea, they have no need to leave the water.
The only genus that has retained the enlarged ventral scales is the sea kraits, Laticauda, with only five species. These snakes are considered to be more primitive, as they still spend much of their time on land, where their ventral scales afford them the necessary grip. Laticauda species are also the only sea snakes with internasal scales; that is, their nostrils are not located dorsally.
Since a snake's tongue can fulfill its olfactory function more easily under water, its action is short compared to that of terrestrial snake species. Only the forked tips protrude from the mouth through a divided notch in the middle of the rostral scale. The nostrils have valves consisting of a specialized spongy tissue to exclude water, and the windpipe can be drawn up to where the short nasal passage opens into the roof of the mouth. This is an important adaptation for an animal that must surface to breathe, but may have its head partially submerged when doing so. The lung has become very large and extends almost the entire length of the body, although the rear portion is thought to have developed to aid buoyancy rather than to exchange gases. The extended lung possibly also serves to store air for dives.
Most species of sea snakes are able to respire through the top of their skin. This is unusual for reptiles, because their skin is thick and scaly, but experiments with the black-and-yellow sea snake, Hydrophis platurus (a pelagic species), have shown this species can satisfy about 25% of its oxygen requirements in this manner, which allows for prolonged dives.
Like other land animals that have adapted to life in a marine environment, sea snakes ingest considerably more salt than their terrestrial relatives through their diets, when seawater is inadvertently swallowed. Because of this, a more effective means of regulating the salt concentration of their blood is required. In sea snakes, the posterior sublingual glands, located under and around the tongue sheath, allow them to expel salt with their tongue action.
Scalation among sea snakes is highly variable. As opposed to terrestrial snake species that have imbricate scales to protect against abrasion, the scales of most pelagic sea snakes do not overlap. Reef-dwelling species, such as Aipysurus, do have imbricate scales to protect against the sharp coral. The scales themselves may be smooth, keeled, spiny, or granular, the latter often looking like warts. Pelamis has body scales that are "peg-like", while those on its tail are juxtaposed hexagonal plates.
Sensory abilities
Vision, chemoreception (tongue-flicking), and hearing are important senses for terrestrial snakes, but these stimuli become distorted in water. The poor visibility, chemical dilution, and limitation of ground-borne vibrations under water suggest that sea snakes and sea kraits may have unique sensory abilities to compensate for the relative lack of other sensory cues.
Relatively little is known about sea snake vision. A study of photoreceptors in the retina of spine-bellied, Lapemis curtus, and horned, Acalyptophis peronii, sea snakes found three classes of opsins all from cone cells. Despite the absence of rod cells in sea snake eyes, Simeos et al. found the rhodopsin (rh1), the opsin of the rods, still expressed suggesting that in sea snakes some cones may be transmuted rods. Behavioural observations indicate that vision has a limited role for catching prey and mate selection, but sound vibrations and chemoreception may be important. One study identified small sensory organs on the head of Lapemis curtus similar to the mechanoreceptors in alligators and aquatic snake Acrochodus that are used to sense the movement of fish prey. Westhoff et al. recorded auditory brain responses to vibration underwater in Lapemis curtus, which are sensitive enough to detect movement in prey, but were not as sensitive as fish lateral line systems. Similarly, vision appears to be of limited importance for finding mates. Shine experimented with applying skin secretions (pheromones) to snake-like objects to see if male turtle-headed sea snakes, Emydocephalus annulatus, are attracted to female pheromones. Shine found that although vision may be useful over short distances (less than ), pheromones are more important once the male comes in physical contact with an object.
The olive sea snake, Aipysurus laevis, has been found to have photoreceptors in the skin of its tail, allowing it to detect light and presumably ensuring it is completely hidden, including its tail, inside coral holes during the day. While other species have not been tested, A. laevis possibly is not unique among sea snakes in this respect.
Other unique senses, such as electromagnetic reception and pressure detection, have been proposed for sea snakes, but scientific studies have yet to be performed to test these senses.
Distribution and habitat
Sea snakes are mostly confined to the warm tropical waters of the Indian Ocean and the western Pacific Ocean, with a few species found well out into Oceania. The geographic range of one species, Pelamis platurus, is wider than that of any other reptile species, except for a few species of sea turtles. It extends from the east coast of Africa, from Djibouti in the north to Cape Town in the south, across the Indian Ocean, the Pacific, south as far as the northern coast of New Zealand, all the way to the western coast of the Americas, where it occurs from northern Peru in the south (including the Galápagos Islands) to the Gulf of California in the north. Isolated specimens have been found as far north as San Diego and Oxnard in the United States.
Sea snakes do not occur in the Atlantic Ocean. Pelamis possibly would be found there were it not for the cold currents off Namibia and western South Africa that keep it from crossing into the eastern South Atlantic, or south of 5°S latitude along the South American west coast. Sea snakes do not occur in the Red Sea, believed to be due to its increased salinity, so no danger exists of them crossing through the Suez Canal. A lack of salinity is also thought to be the reason why Pelamis has not crossed into the Caribbean via the Panama Canal.
Despite their marine adaptations, most sea snakes prefer shallow waters near land, around islands, and especially somewhat sheltered waters, as well as near estuaries. They may swim up rivers and have been reported as far as from the sea. Others, such as P. platurus, are pelagic and are found in drift lines, slicks of floating debris brought together by surface currents. Some sea snakes inhabit mangrove swamps and similar brackishwater habitats, and two landlocked freshwater forms are found: Hydrophis semperi occurs in Lake Taal in the Philippines, and Laticauda crockeri in Lake Tegano on Rennell Island in the Solomon Islands.
Behavior
Sea snakes are generally reluctant to bite, and are usually considered to be mild-tempered, although variation is seen among species and individuals. Some species, such as P. platurus, which feed by simply gulping down their prey, are more likely to bite when provoked because they seem to use their venom more for defense. Others, such as Laticauda spp., use their venom for prey immobilization. Sea snakes are often handled without concern by local fishermen who unravel and toss them back into the water barehanded, usually without getting bitten, when the snakes frequently become entangled in fishing nets. Species reported as much more aggressive include Aipysurus laevis, Astrotia stokesii, Enhydrina schistosa, Enhydrina zweifeli, and Hydrophis ornatus.
On land, their movements become very erratic. They crawl awkwardly in these situations and can become quite aggressive, striking wildly at anything that moves, although they are unable to coil and strike in the manner of terrestrial snakes.
Sea snakes appear to be active both day and night. In the morning, and sometimes late in the afternoon, they can be seen at the surface basking in the sunlight, and they dive when disturbed. They have been reported swimming at depths over , and can remain submerged for as long as a few hours, possibly depending on temperature and degree of activity.
Sea snakes have been sighted in huge numbers. For example, in 1932, a steamer in the Strait of Malacca, off the coast of Malaysia, reported sighting "millions" of Astrotia stokesii, a relative of Pelamis; these reportedly formed a line of snakes wide and long. The cause of this phenomenon is unknown, although it likely has to do with reproduction. They can sometimes be seen swimming in schools of several hundred, and many dead specimens have been found on beaches after typhoons.
Ecology
They feed on small fish and occasionally young octopus. They are often associated with the sea snake barnacle (Platylepas ophiophila), which attaches to their skin.
Reproduction
Except for a single genus, all sea snakes are ovoviviparous; the young are born alive in the water where they live their entire lives. In some species, the young are quite large, up to half as long as the mother. The one exception is the genus Laticauda, which is oviparous; its five species all lay their eggs on land.
Venom
Like their relatives in the family Elapidae, the majority of sea snakes are highly venomous. They rarely inject their venom when biting, so venomous bites to humans are rare. For example, Hydrophis platurus has a venom more potent than any terrestrial snake species in Costa Rica based on LD50, but despite its abundance in the waters off its western coast, few human fatalities have been reported. The death of a trawler fisherman in Australian waters during 2018 was reported to be the region's first sea snake fatality since a pearl diver was killed in 1935.
Bites in which envenomation does occur are usually painless and may not even be noticed when contact is made. Teeth may remain in the wound. Usually, little or no swelling occurs, and rarely are any nearby lymph nodes affected. The most important symptoms are rhabdomyolysis (rapid breakdown of skeletal muscle tissue) and paralysis. Early symptoms include headache, a thick-feeling tongue, thirst, sweating, and vomiting. The venom is very slow acting and symptoms that happen from little as 30 minutes to several hours after the bite include generalized aching, stiffness, and tenderness of muscles all over the body. Passive stretching of the muscles is also painful, and trismus, which is similar to tetanus, is common. This is followed later on by symptoms typical of other elapid envenomations, a progressive flaccid paralysis, starting with ptosis and paralysis of voluntary muscles. Paralysis of muscles involved in swallowing and respiration can be fatal.
Vick et al (1975) estimated that the LD50 of three sea snake venoms (H. platurus, L. semifasciata and L. laticaudata) for a 70 kg human range from 7.7 to 21 mg. Data from the only sea snake venom conducted in monkeys at that time suggested that primates were slightly more resistant to the venom effects on a dose response basis than mice. Ishikawa et al (1985) indicated a substantially lower binding affinity between sea snake neurotoxin and human and chimpanzee AChR's compared to that in other animals. In humans, the venom targets appear mainly to be the cell walls of voluntary (skeletal) muscles and distal tubular portions of the kidney including the loop of Henle, the second convoluted tubule and the collecting tubules. Sitprija et al (1973) found evidence of tubular necrosis throughout all portions of the renal tubules in two patients severely envenomated by sea snakes. Sea snake venoms in humans are thus more often myotoxic and/or nephrotoxic rather than neurotoxic.
Taxonomy
Sea snakes were at first regarded as a unified and separate family, the Hydrophiidae, that later came to comprise two subfamilies: the Hydrophiinae, or true/aquatic sea snakes (now 6 genera with 64 species), and the more primitive Laticaudinae, or sea kraits (one genus, Laticauda, with eight species). Eventually, as just how closely related the sea snakes are to the elapids became clear, the taxonomic situation became less well-defined. Some taxonomists responded by moving the sea snakes to the Elapidae. Most taxonomists now place the sea snakes in the elapid subfamilies Hydrophiinae and Laticaudinae, although the latter may be omitted if Laticauda is included in the Hydrophiinae. Unlike the traditional Hydrophiinae, the Hydrophiinae as currently seen also includes Australasian terrestrial elapids.
Molecular studies
Molecular data studies suggest all three monotypic semiaquatic genera (Ephalophis, Parahydrophis and Hydrelaps) are early diverging lineages.
Captivity
At best, sea snakes make difficult captives. Ditmars (1933) described them as nervous and delicate captives that usually refuse to eat, preferring only to hide in the darkest corner of the tank. Over 50 years later, Mehrtens wrote in 1987 that although they were rarely displayed in Western zoological parks, some species were regularly on display in Japanese aquariums. The available food supply limits the number of species that can be kept in captivity, since some have diets that are too specialized. Also, some species appear intolerant of handling, or even being removed from the water. Regarding their requirements in captivity, the Laticauda species need to be able to exit the water somewhere at about , along with a submerged shelter. Species that have done relatively well in captivity include the ringed sea snake, Hydrophis cyanocinctus, which feeds on fish and eels in particular. Pelamis platurus has done especially well in captivity, accepting small fish, including goldfish. Housing them in round tanks, or in rectangular tanks with well-rounded corners, prevents snakes from damaging their snouts on the sides.
Conservation status
Most sea snakes are not on the CITES protection lists. One species, Laticauda crockeri, is classified as vulnerable. Several species of Aipysurus are listed with conservation status of greater concern, the Timor species A. fuscus is known to be endangered, and two others found in seas north of Australia, the leaf-scaled A. foliosquama and short-nosed A. apraefrontalis, are classified as critically endangered according to the IUCN Red List of Threatened Species.
| Biology and health sciences | Reptiles | null |
27401902 | https://en.wikipedia.org/wiki/Nymphaea%20thermarum | Nymphaea thermarum | Nymphaea thermarum, also known as Pygmy Rwandan water lily, is a species of water lily that is endemic to Rwanda. Once thought to be extinct in the wild, all wild plants were believed to be lost due to destruction of its native habitat, but it was thought to be saved from extinction when it was grown from seed at the Royal Botanic Gardens, Kew in 2009. A previously-unknown wild population was discovered in 2023.
Description
Vegetative characteristics
It is a diminutive, aquatic, rhizomatous herb with 1–2(–5) cm long rhizomes. The peltate, petiolate, glabrous, orbicular to suborbicular leaves have a 2.8–3.2 cm long, and 2.5–3 cm wide lamina. The lobes of the lamina are overlapping each other, or are almost parallel. The petiole is 4–6(–8) cm long.
Generative characteristics
The up to 2 cm wide, hermaphroditic, incompletely protogynous, white flowers have 1.5–3 cm long peduncles. The four green, lanceolate sepals with a round apex are 1.7–1.8 cm long, and 0.6–0.7 cm wide. The 6–8 white petals are 1.5–1.6 cm long, and 4 mm wide. The androecium consists of up to 16 stamens with a sterile apical appendage and they gradually decrease in size from the 9–10 mm long outer stamens to the 5–6 mm long inner stamens. The filament is up to 5–6 mm long, likewise, the connective is 5–6 mm long, and the anthers are 1.5–4 mm long. The gynoecium consists of 7–9 basally fused, 4–5 mm long carpels with a 2mm long, and 1 mm wide stigma forming a stigmatic disk. The up to 1.2–1.5 cm wide fruit bears hundreds of bright red to brown, arillate seeds.
Cytology
The diploid chromosome count is 2n = 28. The genome size is 498.78 Mb. The chloroplast genomes of Nymphaea thermarum and Nymphaea heudelotii are identical.
Taxonomy
Nymphaea thermarum was published by the German botanist Eberhard Fischer in 1988. The type specimen was collected by Fischer in hot springs South of Nyakabuye, Rwanda on the 25th of April 1987. Is is placed within the subgenus Nymphaea subg. Brachyceras.
Etymology
The specific epithet thermarum refers to the hot spring and temperature that provided its native habitat.
Breeding
They can self-pollinate, and after blooming the flower stalk bends so the fruit contacts the mud. The fruit contains 300 to 400 seeds. The sepals are slightly hairy, and as large as the flower's petals. The plant is a tropical day bloomer displaying protogynous flowering patterns, opening early in the morning on the first day with female floral functioning, closing in early afternoon, and opening on the second day with male functionality. It is in the Nymphaea subgenus Brachyceras, though the leaves are more typical of the subgenus Nymphaea. It apparently does not form tubers. Seeds are large for plants in subgenus Brachyceras. The lifespan of Nymphaea thermarum can be greater than 10 years.
Conservation
The plant's native habitat was damp mud formed by the overflow of a freshwater hot spring in Mashyuza, southwest Rwanda. It was thought to have become extinct in the wild around 2008, when local farmers began using the spring for agriculture. The farmers cut off the flow of the spring, which dried up the tiny area—just a few square metres—that was believed to be the entire habitat. Before the extinction of the first known population, Fischer sent some specimens to Bonn Botanic Gardens in Germany when he saw that their habitat was fragile. The plants were kept alive at the gardens, but botanists could not solve the problem of propagating them from seed.
The first published occurrence of N. thermarum germination was by Carlos Magdalena, at the Royal Botanical Gardens, Kew. By placing the seeds and seedlings into pots of loam surrounded by water of the same level in a environment, eight began to flourish and mature within weeks and in November 2009, the waterlilies flowered for the first time. During this time, a rat had eaten one of the last two cultivated plants in Germany. With the germination problem solved, Magdalena says that the tiny plants are easy to grow, giving it potential to be grown as a houseplant. In January 2014, a surviving water lily was stolen from the Royal Botanic Gardens. Botanic gardens have been criticised for not providing plant material to repatriate to Rwanda.
Uses
It has been proposed to be used as a model species for basal angiosperms, due to its small size, rapid lifecycle, and small genome. For instance, together with Nymphaea dimorpha it has been used to study seed evolution.
| Biology and health sciences | Nymphaeales | Plants |
27404990 | https://en.wikipedia.org/wiki/Haskell | Haskell | Haskell () is a general-purpose, statically-typed, purely functional programming language with type inference and lazy evaluation. Designed for teaching, research, and industrial applications, Haskell has pioneered several programming language features such as type classes, which enable type-safe operator overloading, and monadic input/output (IO). It is named after logician Haskell Curry. Haskell's main implementation is the Glasgow Haskell Compiler (GHC).
Haskell's semantics are historically based on those of the Miranda programming language, which served to focus the efforts of the initial Haskell working group. The last formal specification of the language was made in July 2010, while the development of GHC continues to expand Haskell via language extensions.
Haskell is used in academia and industry. , Haskell was the 28th most popular programming language by Google searches for tutorials, and made up less than 1% of active users on the GitHub source code repository.
History
After the release of Miranda by Research Software Ltd. in 1985, interest in lazy functional languages grew. By 1987, more than a dozen non-strict, purely functional programming languages existed. Miranda was the most widely used, but it was proprietary software. At the conference on Functional Programming Languages and Computer Architecture (FPCA '87) in Portland, Oregon, there was a strong consensus that a committee be formed to define an open standard for such languages. The committee's purpose was to consolidate existing functional languages into a common one to serve as a basis for future research in functional-language design.
Haskell 1.0 to 1.4
Haskell was developed by a committee, attempting to bring together off the shelf solutions where possible.
Type classes, which enable type-safe operator overloading, were first proposed by Philip Wadler and Stephen Blott to address the ad-hoc handling of equality types and arithmetic overloading in languages at the time.
In early versions of Haskell up until and including version 1.2, user interaction and input/output (IO) were handled by both streams based and continuation based mechanisms which were widely considered unsatisfactory. In version 1.3, monadic IO was introduced, along with the generalisation of type classes to higher kinds (type constructors). Along with "do notation", which provides syntactic sugar for the Monad type class, this gave Haskell an effect system that maintained referential transparency and was convenient.
Other notable changes in early versions were the approach to the 'seq' function, which creates a data dependency between values, and is used in lazy languages to avoid excessive memory consumption; with it moving from a type class to a standard function to make refactoring more practical.
The first version of Haskell ("Haskell 1.0") was defined in 1990. The committee's efforts resulted in a series of language definitions (1.0, 1.1, 1.2, 1.3, 1.4).
Haskell 98
In late 1997, the series culminated in Haskell 98, intended to specify a stable, minimal, portable version of the language and an accompanying standard library for teaching, and as a base for future extensions. The committee expressly welcomed creating extensions and variants of Haskell 98 via adding and incorporating experimental features.
In February 1999, the Haskell 98 language standard was originally published as The Haskell 98 Report. In January 2003, a revised version was published as Haskell 98 Language and Libraries: The Revised Report. The language continues to evolve rapidly, with the Glasgow Haskell Compiler (GHC) implementation representing the current de facto standard.
Haskell 2010
In early 2006, the process of defining a successor to the Haskell 98 standard, informally named Haskell Prime, began. This was intended to be an ongoing incremental process to revise the language definition, producing a new revision up to once per year. The first revision, named Haskell 2010, was announced in November 2009 and published in July 2010.
Haskell 2010 is an incremental update to the language, mostly incorporating several well-used and uncontroversial features previously enabled via compiler-specific flags.
Hierarchical module names. Module names are allowed to consist of dot-separated sequences of capitalized identifiers, rather than only one such identifier. This lets modules be named in a hierarchical manner (e.g., Data.List instead of List), although technically modules are still in a single monolithic namespace. This extension was specified in an addendum to Haskell 98 and was in practice universally used.
The foreign function interface (FFI) allows bindings to other programming languages. Only bindings to C are specified in the Report, but the design allows for other language bindings. To support this, data type declarations were permitted to contain no constructors, enabling robust nonce types for foreign data that could not be constructed in Haskell. This extension was also previously specified in an Addendum to the Haskell 98 Report and widely used.
So-called n+k patterns (definitions of the form fact (n+1) = (n+1) * fact n) were no longer allowed. This syntactic sugar had misleading semantics, in which the code looked like it used the (+) operator, but in fact desugared to code using (-) and (>=).
The rules of type inference were relaxed to allow more programs to type check.
Some syntax issues (changes in the formal grammar) were fixed: pattern guards were added, allowing pattern matching within guards; resolution of operator fixity was specified in a simpler way that reflected actual practice; an edge case in the interaction of the language's lexical syntax of operators and comments was addressed, and the interaction of do-notation and if-then-else was tweaked to eliminate unexpected syntax errors.
The LANGUAGE pragma was specified. By 2010, dozens of extensions to the language were in wide use, and GHC (among other compilers) provided the LANGUAGE pragma to specify individual extensions with a list of identifiers. Haskell 2010 compilers are required to support the Haskell2010 extension and are encouraged to support several others, which correspond to extensions added in Haskell 2010.
Future standards
The next formal specification had been planned for 2020. On 29 October 2021, with GHC version 9.2.1, the GHC2021 extension was released. While this is not a formal language spec, it combines several stable, widely-used GHC extensions to Haskell 2010.
Features
Haskell features lazy evaluation, lambda expressions, pattern matching, list comprehension, type classes and type polymorphism. It is a purely functional programming language, which means that functions generally have no side effects. A distinct construct exists to represent side effects, orthogonal to the type of functions. A pure function can return a side effect that is subsequently executed, modeling the impure functions of other languages.
Haskell has a strong, static type system based on Hindley–Milner type inference. Its principal innovation in this area is type classes, originally conceived as a principled way to add overloading to the language, but since finding many more uses.
The construct that represents side effects is an example of a monad: a general framework which can model various computations such as error handling, nondeterminism, parsing and software transactional memory. They are defined as ordinary datatypes, but Haskell provides some syntactic sugar for their use.
Haskell has an open, published specification, and multiple implementations exist. Its main implementation, the Glasgow Haskell Compiler (GHC), is both an interpreter and native-code compiler that runs on most platforms. GHC is noted for its rich type system incorporating recent innovations such as generalized algebraic data types and type families. The Computer Language Benchmarks Game also highlights its high-performance implementation of concurrency and parallelism.
An active, growing community exists around the language, and more than 5,400 third-party open-source libraries and tools are available in the online package repository Hackage.
Code examples
A "Hello, World!" program in Haskell (only the last line is strictly necessary):
module Main (main) where -- not needed in interpreter, is the default in a module file
main :: IO () -- the compiler can infer this type definition
main = putStrLn "Hello, World!"
The factorial function in Haskell, defined in a few different ways (the first line is the type annotation, which is optional and is the same for each implementation):
factorial :: (Integral a) => a -> a
-- Using recursion (with the "ifthenelse" expression)
factorial n = if n < 2
then 1
else n * factorial (n - 1)
-- Using recursion (with pattern matching)
factorial 0 = 1
factorial n = n * factorial (n - 1)
-- Using recursion (with guards)
factorial n
| n < 2 = 1
| otherwise = n * factorial (n - 1)
-- Using a list and the "product" function
factorial n = product [1..n]
-- Using fold (implements "product")
factorial n = foldl (*) 1 [1..n]
-- Point-free style
factorial = foldr (*) 1 . enumFromTo 1
Using Haskell's Fixed-point combinator allows this function to be written without any explicit recursion.
import Data.Function (fix)
factorial = fix fac
where fac f x
| x < 2 = 1
| otherwise = x * f (x - 1)
As the Integer type has arbitrary-precision, this code will compute values such as factorial 100000 (a 456,574-digit number), with no loss of precision.
An implementation of an algorithm similar to quick sort over lists, where the first element is taken as the pivot:
-- Type annotation (optional, same for each implementation)
quickSort :: Ord a => [a] -> [a]
-- Using list comprehensions
quickSort [] = [] -- The empty list is already sorted
quickSort (x:xs) = quickSort [a | a <- xs, a < x] -- Sort the left part of the list
++ [x] ++ -- Insert pivot between two sorted parts
quickSort [a | a <- xs, a >= x] -- Sort the right part of the list
-- Using filter
quickSort [] = []
quickSort (x:xs) = quickSort (filter (<x) xs)
++ [x] ++
quickSort (filter (>=x) xs)
Implementations
All listed implementations are distributed under open source licenses.
Implementations that fully or nearly comply with the Haskell 98 standard include:
The Glasgow Haskell Compiler (GHC) compiles to native code on many different processor architectures, and to ANSI C, via one of two intermediate languages: C--, or in more recent versions, LLVM (formerly Low Level Virtual Machine) bitcode. GHC has become the de facto standard Haskell dialect. There are libraries (e.g., bindings to OpenGL) that work only with GHC. GHC was also distributed with the Haskell platform.
Jhc, a Haskell compiler written by John Meacham, emphasizes speed and efficiency of generated programs and exploring new program transformations.
Ajhc is a fork of Jhc.
The Utrecht Haskell Compiler (UHC) is a Haskell implementation from Utrecht University. It supports almost all Haskell 98 features plus many experimental extensions. It is implemented using attribute grammars and is primarily used for research on generated type systems and language extensions.
Implementations no longer actively maintained include:
The Haskell User's Gofer System (Hugs) is a bytecode interpreter. It was once one of the implementations used most widely, alongside the GHC compiler, but has now been mostly replaced by GHCi. It also comes with a graphics library.
HBC is an early implementation supporting Haskell 1.4. It was implemented by Lennart Augustsson in, and based on, Lazy ML. It has not been actively developed for some time.
nhc98 is a bytecode compiler focusing on minimizing memory use.
The York Haskell Compiler (Yhc) was a fork of nhc98, with the goals of being simpler, more portable and efficient, and integrating support for Hat, the Haskell tracer. It also had a JavaScript backend, allowing users to run Haskell programs in web browsers.
Implementations not fully Haskell 98 compliant, and using a variant Haskell language, include:
Eta and Frege are dialects of Haskell targeting the Java virtual machine.
Gofer is an educational dialect of Haskell, with a feature called constructor classes, developed by Mark Jones. It is supplanted by Haskell User's Gofer System (Hugs).
Helium, a newer dialect of Haskell. The focus is on making learning easier via clearer error messages by disabling type classes as a default.
Notable applications
Agda is a proof assistant written in Haskell.
Cabal is a tool for building and packaging Haskell libraries and programs.
Darcs is a revision control system written in Haskell, with several innovative features, such as more precise control of patches to apply.
Glasgow Haskell Compiler (GHC) is also often a testbed for advanced functional programming features and optimizations in other programming languages.
Git-annex is a tool to manage (big) data files under Git version control. It also provides a distributed file synchronization system (git-annex assistant).
Linspire Linux chose Haskell for system tools development.
Pandoc is a tool to convert one markup format into another.
Pugs is a compiler and interpreter for the programming language then named Perl 6, but since renamed Raku.
TidalCycles is a domain special language for live coding musical patterns, embedded in Haskell.
Xmonad is a window manager for the X Window System, written fully in Haskell.
GarganText is a collaborative tool to map through semantic analysis texts on any web browser, written fully in Haskell and PureScript, which is used for instance in the research community to draw up state-of-the-art reports and roadmaps.
Industry
Bluespec SystemVerilog (BSV) is a language extension of Haskell, for designing electronics. It is an example of a domain-specific language embedded into Haskell. Further, Bluespec, Inc.'s tools are implemented in Haskell.
Cryptol, a language and toolchain for developing and verifying cryptography algorithms, is implemented in Haskell.
Facebook implements its anti-spam programs in Haskell, maintaining the underlying data access library as open-source software.
The Cardano blockchain platform is implemented in Haskell.
GitHub implemented Semantic, an open-source library for analysis, diffing, and interpretation of untrusted source code, in Haskell.
Standard Chartered's financial modelling language Mu is syntactic Haskell running on a strict runtime.
seL4, the first formally verified microkernel, used Haskell as a prototyping language for the OS developer. At the same time, the Haskell code defined an executable specification with which to reason, for automatic translation by the theorem-proving tool. The Haskell code thus served as an intermediate prototype before final C refinement.
Target stores' supply chain optimization software is written in Haskell.
Co–Star
Mercury Technologies' back end is written in Haskell.
Web
Notable web frameworks written for Haskell include:
IHP
Servant
Snap
Yesod
Criticism
Jan-Willem Maessen, in 2002, and Simon Peyton Jones, in 2003, discussed problems associated with lazy evaluation while also acknowledging the theoretical motives for it. In addition to purely practical considerations such as improved performance, they note that lazy evaluation makes it more difficult for programmers to reason about the performance of their code (particularly its space use).
Bastiaan Heeren, Daan Leijen, and Arjan van IJzendoorn in 2003 also observed some stumbling blocks for Haskell learners: "The subtle syntax and sophisticated type system of Haskell are a double edged sword—highly appreciated by experienced programmers but also a source of frustration among beginners, since the generality of Haskell often leads to cryptic error messages." To address the error messages researchers from Utrecht University developed an advanced interpreter called Helium, which improved the user-friendliness of error messages by limiting the generality of some Haskell features. In particular it disables type classes by default.
Ben Lippmeier designed Disciple as a strict-by-default (lazy by explicit annotation) dialect of Haskell with a type-and-effect system, to address Haskell's difficulties in reasoning about lazy evaluation and in using traditional data structures such as mutable arrays. He argues (p. 20) that "destructive update furnishes the programmer with two important and powerful tools ... a set of efficient array-like data structures for managing collections of objects, and ... the ability to broadcast a new value to all parts of a program with minimal burden on the programmer."
Robert Harper, one of the authors of Standard ML, has given his reasons for not using Haskell to teach introductory programming. Among these are the difficulty of reasoning about resource use with non-strict evaluation, that lazy evaluation complicates the definition of datatypes and inductive reasoning, and the "inferiority" of Haskell's (old) class system compared to ML's module system.
Haskell's build tool, Cabal, has historically been criticized for poorly handling multiple versions of the same library, a problem known as "Cabal hell". The Stackage server and Stack build tool were made in response to these criticisms. Cabal itself has since addressed this problem by borrowing ideas from Nix,, with the new approach becoming the default in 2019.
Related languages
Clean is a close, slightly older relative of Haskell. Its biggest deviation from Haskell is in the use of uniqueness types instead of monads for input/output (I/O) and side effects.
A series of languages inspired by Haskell, but with different type systems, have been developed, including:
Agda, a functional language with dependent types.
Cayenne, with dependent types.
Elm, a functional language to create web front-end apps, no support for user-defined or higher-kinded type classes or instances.
Epigram, a functional language with dependent types suitable for proving properties of programs.
Idris, a general purpose functional language with dependent types, developed at the University of St Andrews.
PureScript transpiles to JavaScript.
Ωmega, a strict language that allows introduction of new kinds, and programming at the type level.
Other related languages include:
Curry, a functional/logic programming language based on Haskell.
Notable Haskell variants include:
Generic Haskell, a version of Haskell with type system support for generic programming.
Hume, a strict functional language for embedded systems based on processes as stateless automata over a sort of tuples of one element mailbox channels where the state is kept by feedback into the mailboxes, and a mapping description from outputs to channels as box wiring, with a Haskell-like expression language and syntax.
Conferences and workshops
The Haskell community meets regularly for research and development activities. The main events are:
International Conference on Functional Programming (ICFP)
Haskell Symposium (formerly the Haskell Workshop)
Haskell Implementors Workshop
Commercial Users of Functional Programming (CUFP)
ZuriHac, kind of Hackathon held every year in Zurich
Starting in 2006, a series of organized hackathons has occurred, the Hac series, aimed at improving the programming language tools and libraries.
| Technology | Programming languages | null |
145550 | https://en.wikipedia.org/wiki/Flechette | Flechette | A flechette or flèchette ( ) is a pointed, fin-stabilized steel projectile. The name comes from French (from flèche), meaning "little arrow" or "dart", and sometimes retains the grave accent in English: flèchette. They have been used as ballistic weapons since World War I. Delivery systems and methods of launching flechettes vary, from a single shot, to thousands in a single explosive round. The use of flechettes as antipersonnel weapons has been controversial.
Air-dropped
The weapons were designed to be dropped from an aircraft. They contained no explosive charge but as they fell they developed significant kinetic energy making them lethal and able to easily penetrate soft cover such as jungle canopy, several inches of sand or light armor.
During World War I, flechettes were dropped from aircraft to attack infantry and were able to pierce helmets.
Later the U.S. used Lazy Dog bombs, which are small, unguided kinetic projectiles typically about in length, in diameter, and weighing about . Lazy Dog munitions were simple and cheap; they could be dropped in huge numbers in a pass. Though their effects were often no less indiscriminate than other projectiles, they did not leave unexploded ordnance (UXO) that could be active years after a conflict ended. Lazy Dog projectiles were used primarily during the Korean and Vietnam Wars.
Small-arms ammunition
The excellent ballistic performance and armor-piercing potential of flechettes have made the development and integration of this class of munition attractive to small-arms manufacturers. A number of attempts have been made to field flechette-firing small arms.
Work at Johns Hopkins University in the 1950s led to the development of the direct injection antipersonnel chemical biological agent (DIACBA), where flechettes were grooved, hollow pointed, or otherwise milled to retain a quantity of chemical or biological warfare agent to be delivered through a ballistic wound. The initial work was with the nerve agent VX, which had to be thickened to deliver a reliable dose. Eventually this was replaced by a particulate carbamate. The US Biological Program also had a microflechette to deliver either botulinum toxin A or saxitoxin, the M1 biodart, which resembled a 7.62 mm rifle cartridge. The USSR had the AO-27 rifle as well as the APS amphibious rifle, and other countries have their own flechette rounds.
A number of prototype flechette-firing weapons were developed as part of the long-running Special Purpose Individual Weapon (SPIW) project. The Steyr-Mannlicher ACR rifle was a prototype flechette-firing assault rifle built for the US Army's Advanced Combat Rifle program of 1989–90.
A variation of the flechette addressing its difficulties is the SCMITR, developed as part of the Close Assault Weapon System, or CAWS, project. Selective-fire shotguns were used to fire flechettes designed to retain the exterior ballistics and penetration of standard flechettes, but increase wounding capacity through a wider wound path.
Shotguns
During the Vietnam War the United States employed 12-gauge combat shotguns using flechette loads. These plastic-cased shells were issued on a limited trial basis during the Vietnam War. Cartridges manufactured by the Western Cartridge Company contained 20 flechettes, each long and weighing ; Federal Cartridge Company rounds contained 25. The flechettes were packed in a plastic cup with granulated white polyethylene to maintain alignment with the bore axis, and supported by a metal disk to prevent penetration of the over-powder wad during acceleration down the bore. The tips of the flechettes were exposed in the Federal cartridges, but concealed by a conventional star crimp in WCC's cartridges. The flechettes demonstrated flatter trajectories over longer ranges than spherical buckshot, but combat effectiveness did not justify continued production.
Rocket and artillery use
Smaller flechettes were used in special artillery shells called "beehive" rounds (so named for the very distinctive whistling buzz made by thousands of flechettes flying downrange at supersonic speeds) and intended for use against troops in the open – a ballistic shell packed with flechettes was fired and set off by a mechanical time fuse, scattering flechettes in an expanding cone.
During the Vietnam War 105 mm howitzer batteries and tanks (90 mm guns) used flechette rounds to defend themselves against massed infantry attacks. The ubiquitous 105 mm M40 recoilless rifle was primarily used as an anti-tank weapon. However, it could also be used in an anti-personnel role with the use of flechette rounds. The widely used Carl Gustaf 8.4 cm recoilless rifle also uses an Area Defence Munition designed as a close-range anti-personnel round. It fires 1,100 flechettes over a wide area. The US Air Force used rockets with WDU-4/A flechette warheads.
The 70 mm Hydra 70 rocket currently in service with the US Armed forces can be fitted with an anti-personnel (APERS) warhead containing 1,179 flechettes. They are carried by attack helicopters such as the AH-64 Apache and the AH-1 Cobra.
Israel-Palestine conflict
Israeli authorities have publicly confirmed the use of flechettes against Palestinians in Gaza since at least 2001. International humanitarian organizations documented numerous instances of flechette shells being used by the IDF during Operation Cast Lead, resulting in the deaths of at least nine civilians. According to Amnesty International, the use of flechette shells in densely populated civilian areas violated the international prohibition on indiscriminate attacks. In 2001, Israeli officials stated that "The Israeli military obtained these weapons from the USA after the 1973 war and we have thousands of old shells in warehouses," suggesting that flechettes may have been used against Palestinians since the 1970s.
Russo-Ukrainian war
Flechettes have been used during the 2022 Russian invasion of Ukraine, where samples of the projectiles were recovered in the mass graves in Bucha. A witness described munitions bursting overhead and littering the area with 3 cm flechettes. A British munitions expert reviewed photographs of the flechettes and concluded that they likely came from a 122 mm 3Sh1 artillery round. A speaker for the Ukrainian Ground Forces stated that Ukraine's military does not use shells with flechettes.
| Technology | Ammunition | null |
145672 | https://en.wikipedia.org/wiki/Manx%20cat | Manx cat | The Manx cat (, in earlier times often spelled Manks) is a breed of domestic cat (Felis catus) originating on the Isle of Man, with a mutation that shortens the tail. Many Manx have a small stub of a tail, but Manx cats are best known as being entirely tailless; this is the most distinguishing characteristic of the breed, along with elongated hind legs and a rounded head. Manx cats come in all coat colours and patterns, though all-white specimens are rare, and the coat range of the original stock was more limited. Long-haired variants are sometimes considered a separate breed, the Cymric cat.
Manx are prized as skilled hunters, and thus have often been sought by farmers with rodent problems, and been a preferred ship's cat breed. They are said to be social, tame and active. An old local term for the cats on their home island is stubbin or rumpy. Manx have been exhibited in cat shows since the 1800s, with the first known breed standard published in 1903.
History
Origin and folklore
Tailless cats, then called (apparently both singular and plural) in colloquial Manx language, were known by the early 19th century as cats from the Isle of Man, hence the name, where they remain a substantial but declining percentage of the local cat population. The taillessness arose as a natural mutation on the island, though folklore persists that tailless domestic cats were brought there by sea. They are descended from mainland stock of obscure origin. Like all house cats, including nearby British and Irish populations, they are ultimately descended from the African wildcat (Felis lybica) and not from native European wildcats (Felis silvestris), of which the island has long been devoid.
The dominant trait of taillessness arises from a spontaneous mutation, the Manx taillessness gene, that eventually became common on the island because of the limited genetic diversity of island biogeography (an example of the founder effect and, at , of the species-area curve).
In the Manx language, the modern name of the breed is , literally 'cat of Mann' (plural or ), or lit. 'bob-tailed cat'. , used as both a masculine and feminine noun, is also encountered as , and depending on the exact construction, it may be lenited as or . The diminutive word is or , 'kitten' (with various plurals). Manx itself was often spelled Manks in English well into the late 1800s.
There are numerous folktales about the Manx cat, all of them of "relatively recent origin"; they are focused entirely on the lack of a tail, and are devoid of religious, philosophical, or mythical aspects found in the traditional Irish–Norse folklore of the native Manx culture, and in legends about cats from other parts of the world.
The name of the promontory Spanish Head on the coast of the island is often thought to have arisen from the local tale of a ship of the Spanish Armada foundering in the area, though there is no evidence to suggest this actually occurred. Folklore has further claimed that a tailless cat swam ashore from said shipwreck, and thus brought the trait to the island. However, tailless cats are not commonly known in Spain, even if such a shipwreck were proven.
Regardless of the genetic and historical reality, there are various fanciful Lamarckian folktales that seek to explain why the Manx has a truncated tail. In one of them, the biblical Noah closed the door of the Ark when it began to rain, and accidentally cut off the tail of the Manx cat who had almost been left behind. Over the years a number of cartoons have appeared on postcards from the Isle of Man showing scenes in which a cat's tail is being run over and severed by a variety of means including a motorcycle, a reference to motorcycle racing being popular on the island, and an update of the Noah story. Because the gene is so dominant and "invades" other breeds when crossed (often without owner knowledge) with the Manx, there was a folk belief that simply being in the proximity of a Manx cat could cause other breeds to somehow produce tailless kittens.
Another genetically impossible account claimed that the Manx was the hybrid offspring of a cat and a rabbit, purporting to explain why it has no or little tail, long hind legs and a sometimes hopping gait. The cat-rabbit halfbreed tale has been further reinforced by the more widespread "cabbit" folktale.
Populations of tailless cats also exist in a few other places in Europe, most notably Cornwall, only from the Isle of Man. A population on the small, isolated Danish peninsula (former island) of Reersø in the Great Belt may be due to the arrival on the island of cats of Manx origin, by ship. Similar cats are also found in Crimea, a near-island peninsula in the Black Sea, though whether they are genetically related to maritime Manx cats or are a coincidentally similar result of insular genetic diversity limitations, like the unrelated Kuril Islands Bobtail, Karelian Bobtail, Japanese Bobtail, and Indonesian Lombok cats, is unknown. The Manx gene may be related to the similarly dominant tail suppression gene of the recent American Bobtail breed, but Manx, Japanese Bobtails and other short-tailed cats are not used in its breeding program, and the mutation seems to have appeared in the breed spontaneously. Possible relation to the Pixie-bob breed, which also ranges from rumpy to fully tailed, is unknown.
Recognition as a breed
Manx cats have been exhibited in cat shows, as a named, distinct breed (and with the modern spelling "Manx"), since the late 1800s. In that era, few shows provided a Manx division, and exhibited specimens were usually entered under the "Any Other Variety" class, where they often could not compete well unless "exceptionally good in size and markings". Early pet breeding and showing expert Charles Henry Lane, himself the owner of a prize-winning rare white rumpy Manx named Lord Luke, published the first known (albeit informal) breed standard for the Manx in his 1903 Rabbits, Cats and Cavies, but noted that already by the time of his writing "if the judge understood the variety" a Manx would be clearly distinguishable from some other tailless cat being exhibited, "as the make of the animal, its movements and its general character are all distinctive." Not all cat experts of the day were favourable toward the breed; in The Cat: Its Points and Management in Health and Disease, Frank Townend Barton wrote in 1908: "There is nothing to recommend the breed, the loss of the tail in no way enhances its beauty."
The Manx was one of the first breeds recognised by the Cat Fanciers' Association (CFA) (the predominant United States–based pedigreed cat registry, founded in 1908), which has records on the breed in North America going back to the 1920s.
Appearance
Tail (or lack thereof)
Although tail suppression (or tail length variety) is not the sole characteristic feature of the breed, the chief defining one of the Manx cat is its absence of a tail to having a tail of long length, or tail of any length between the two extremes. This is a cat body-type mutation of the spine, caused by a dominant gene. As with the sometimes-tail-suppressed Schipperke dog and Old English Sheepdog, tail suppression does not "breed true" in Manx cats. Attempting to force the tailless trait to breed true by continually breeding tailless Manx cats to tailless Manx cats has led to increased negative, even fatal genetic disorders (see below). Tail length is random throughout a litter of kittens. Manx to non-Manx breeding will usually produce some Manx-type tail varieties in kittens. Whether the shorter tailed kittens of these varieties are labeled Manx is up to the breed standard consulted. Manx cats' tails are classified according to proportional tail length as kittens (the proportion does not change after birth):
Rumpy (rumpie) or dimple rumpy – having no tail at all, though often a small tuft of hair where the tail would have grown from the rump
Riser or rumpy riser – having a bump of cartilage under the fur, most noticeable when the animal is happy and raising its tail end
Stumpy (stumpie) – having a partial tail of vestigial, fused vertebrae, up to about long
Stubby (stubbie), shorty, or short-tailed – having a short tail of non-fused bones, up to about half an average cat tail
Longy (longie), tailed, or taily (tailie) – having a half- to normal-length tail.
Since the early days of breed recognition in the late 19th century, Manx show cats have been rumpy through stumpy specimens, with stubby and longy Manx not qualifying to be shown except in the "Any Other Variety" or household pet class. Kittens with complete tails may be born in a purebred Manx litter, having not inherited the taillessness appearance at all. Depending on the country and cat organization referenced, rumpy, rumpy risers and stumpies are the only Manx cat tail types that fit the breed standard for Manx cats. The longer cat tail lengths seen in some Manx cats are considered a breed fault, although they occur as naturally in the breed, but not as often, as the shorter tails. Although these longer tail types are of purebred Manx ancestry, they do not possess the dominant gene so cannot pass it on. However, since the Manx tail mutation gene is dominant, these longer-tailed purebred Manx cats may still be used in breeding programs and may even be considered in an effort to help avoid the fatal spinal deformities that sometimes result in tailless Manx cats.
The Manx breed is genetically distinct from the Japanese Bobtail breed, another naturally occurring insular breed. The Japanese Bobtail always has at least some tail, ranging from a small "pom" to a stubby but distinct tail, which is kinked or curled and usually has a slightly bulbous and fluffy appearance; by contrast, the Manx has a straight tail when one is present at all. The Japanese Bobtail has a markedly different appearance from the Manx, and is characterized by almond-shaped eyes, a triangular face, long ears, and lean body, like many other Asian breeds. The gene responsible for the bobbed or kinked tail in that breed is recessive and unrelated to the dominant Manx tail-suppression gene; the bobtail gene is not connected to any serious deformities, while the tail-suppression gene can, under certain conditions, give rise to a pattern of sometimes lethal health problems. The Pixie-bob breed also has a short tail, and may be genetically related to the Manx. More will be clear about tail genetics as more genetic studies are done on cat populations and as DNA testing improves; most domestic animal genetic work has been done with dogs and livestock breeds.
Manx (and other tail-suppressed breeds) do not exhibit problems with balance; balance is controlled primarily by the inner ear. In cats, dogs and other large-bodied mammals, balance involves but is not dependent upon the tail (contrast with rats, for whom the tail is a quite significant portion of their body mass).
Since Manx kittens are naturally born with any tail length, from none to long, it was formerly common to surgically dock the longer tails a few days after birth. Although illegal in many jurisdictions (including much of Europe), the practice was formerly recommended, although with the caveat that the commonness of the practice meant that many spurious Manx cats – i.e., random British cats – were altered to resemble the Manx, to defraud unwary buyers.
Body and legs
Manx are medium-sized cats, broad-chested with sloping shoulders and flat sides, and in show condition are firmly muscular and lean, neither bulky nor fatty. Lane reported the original, native breed as ranging typically from ten to twelve pounds for males and eight to ten pounds for females, with many smaller examples but only rare ones larger. The hind legs of Manx are notably longer than the fore legs, causing the rump to be higher than the shoulder and creating a continuous arch from shoulders to rump giving the cat an overall rounded or humped appearance, though the breed is comparatively long when stretched out. The fore legs are strong and straight. The shape is often described as rabbit-like.
Head
Manx cats' heads are rounded in shape, and medium in depth with a long neck. The upright, round-tipped and front-facing ears are largish. The eyes are large, rounded, and prominent, with their outer corners higher than the inner ones. Absent any bloodlines with a dominant alternative eye color (such as blue in Siamese or related ancestry), Manx often have some hue variant of gold eyes, and for show purposes follow the eye colour standards of the same coat colour/pattern in non-Manx short-hairs.
Coat
Manx cats exhibit two coat lengths. Short- or long-haired, all Manx have a thick, double-layered coat. The colour and pattern ranges exhibited should conform to the standards for that type of coat in non-Manx.
The more common short-haired Manx – the original breed – has a coat with a dense, soft, under layer and a longer, coarse outer layer with guard hairs. The overall appearance of the coat is fine, short and lying close to the skin, versus fluffy or voluminous.
The long-haired Manx, known to some cat registries as the Cymric, has a silky-textured double coat of medium length, with "breeches" (i.e. a distinct jump in fur length at the hocks giving the appearance of old-fashioned, baggy, knee-length pants ) belly ruff and neck ruff, tufts of fur between the toes and full "ear furnishings" (hairs in ears). The CFA considers the Cymric to be a variety of Manx and judges it in the short-hair division even though it is long-haired, while The International Cat Association (TICA) judges it in the long-hair division as a distinct Cymric breed. The long-haired variety is of comparatively recent development. Lane wrote in 1903 that the Manx "to the best of my knowledge, information and belief, does not include any long-haired specimens", in his detailed chapter on the breed.
Regardless of coat length, the colours and coat patterns occurring in the breed today run the gamut of virtually all breeds due to extensive cross-breeding, though not all registries may accept all coats as qualifying for breeding or show. The most common coats are tabby, tortoiseshell, calico and solid colours. Widely divergent Manx specimens, including even a colour-point, blue-eyed, long-haired variant of evident Himalayan ancestry, have been celebrated on Isle of Man postage stamps since the 1980s, and recent publications often show marbled and spotted varieties. The original insular stock, however, were of less widespread variation. Lane, having "seen a great many of them" wrote of Manx cats that "[i]t is curious that the colours in this variety seem somewhat limited" and that the breed "does not comprise all the colours usually associated with other short-haired varieties". He reported only very common orange, common orange and white, common cream tabby, uncommon tortoiseshell, and very rare all-white specimens in 1903. Calico and point-coloured are notably absent from this list, as are even today's common colourful tabbies. However, writing in England only five years later, Barton suggested that "the Manx may be of any colour, but probably orange is the most frequently met with."
Specific registries have particular, and differing, standards of points with regard to coloration and patterning. For example, the Governing Council of the Cat Fancy (GCCF) classifies the Manx as a variant of the British Shorthair (BSH), and thus requires that Manx cats to have one of the coat patterns that would be permissible in the BSH rather than any that is exclusive to a "foreign" type (e.g. point colouration). New Zealand Cat Fancy (NZCF) does likewise for colour and markings, but requires a double-coat and other Manx-specific features that GCCF does not. Some other registries are even more restrictive, while others are more liberal.
Variants (sub-breeds)
Four new, consistent varieties have been developed from the Manx (the original version of which is now sometimes consequently called the Shorthair Manx). These are the Cymric (Longhair Manx), the Isle of Man Shorthair and Isle of Man Longhair, and the Tasman Manx, though only the Cymric has garnered widespread acceptance in breed registries .
Cymric (Manx Longhair)
The Cymric or Manx Longhair is a tailless or partially tailed cat of Manx stock, with semi-long to long hair, e.g. as the result of cross-breeding with Himalayan, Persian and other longer-haired breeds early in its development. While its name refers to Wales (), the breed was actually developed in Canada, which has honoured the breed with a commemorative 50-cent coin in 1999.
Simply covering it in their Manx breed standards, the US-based Cat Fanciers' Association (CFA), the Co-ordinating Cat Council of Australia (CCCA), and the UK's Governing Council of the Cat Fancy (GCCF) recognise the variety as a longer-haired Manx rather than "Cymric" (the CFA and CCCA call it the Manx Longhair, while GCCF uses the term Semi-longhair Manx Variant). The majority of cat registries have explicit Cymric standards (published separately or along with Manx). Of the major registries, only the Feline Federation Europe (FFE) does not recognise the breed or sub-breed at all, under any name, (their Manx standard was last updated 17 May 2004).
Isle of Man Shorthair (tailed)
Resembling the British Shorthair, the Isle of Man Shorthair is essentially a fully tailed Manx cat. That is, it is a cat of Manx stock, with Manx features, but without any expression of the Manx taillessness gene. , it is only recognised by New Zealand Cat Fancy (NZCF) with its own breed standard. Any coat colour and pattern acceptable in the British Shorthair is permissible in the IoM Shorthair (the same restriction is applied to the Manx in the NZCF standard), and it requires the double coat of the Manx. In other international registries (e.g. GCCF, who also treat Manx as a British Shorthair variant), such cats are designated "Tailed Manx" and only recognised as Manx breeding stock (they are important as such, since breeding two tailless Manx together results in birth defects), and cannot be show cats.
Isle of Man Longhair (tailed)
Essentially a fully tailed Cymric cat, i.e., a cat of Cymric (and thus Manx) stock, the Isle of Man Longhair has Cymric features, but without expression of the Manx taillessness gene. , it is only recognised as a separate breed by NZCF with a breed standard. Coat colours are limited to those acceptable in the British Shorthair, and requires the double, thick, long coat of the Cymric.
Tasman Manx (curly-coated)
Named after Tasman Sea between Australia and New Zealand, the Tasman Manx is a tailless or partially tailed Manx cat with a curly-haired coat not unlike that of a Selkirk Rex, due a recessive mutation which arose in Manx litters in both Australia and New Zealand. , the breed is only recognised by the NZCF and the Catz Inc. registry (also of New Zealand) with breed standards. The coat may be short or semi-long.
The type arose possibly without existing rex mutation bloodlines (and none of the rex breeds are permitted as out-cross partners with Tasman Manx in Catz breeding guidelines). Depending on length of tail (if any) and coat, kittens may sometimes be termed "Tasman Cymric", "Tasman Isle of Man Shorthair" or "Tasman Isle of Man Longhair", but these are not considered separate breeds. The term "Tasman Rex" has been applied to cats with this gene that do not fall into one of the previously mentioned labels (lacking the Manx face and body shape to qualify), though relation if any to extant Rex mutation breeds is unclear. All of these additional terms beyond "Tasman Manx" appear to be "recognised", even promulgated by NZCF but without breed standards, and even the permissive Catz registry does not include them .
Health and genetics
The Manx taillessness gene is dominant and highly penetrant; kittens from two Manx parents are generally born without any tail. Being homozygous for (having two copies of) the gene is usually lethal in utero, resulting in miscarriage. Thus, tailless cats can carry only one copy of the gene. Because of the danger of having two copies of the taillessness gene, breeders avoid breeding two entirely tailless Manx cats together. Because neither parent carries the tailless allele, a fully tailed Manx bred to another fully tailed Manx results in all fully tailed kittens.
Some partial tails are prone to a form of arthritis that causes the cat severe pain, and in rare cases Manx-bred kittens are born with kinked short tails because of incomplete growth of the tail during development. Stumpy to long tails are sometimes docked at birth as a preventative measure.
"Manx syndrome" or "Manxness" is a colloquial name given to the condition which results when the tailless gene shortens the spine too much. It can seriously damage the spinal cord and the nerves, causing a form of spina bifida, as well as problems with the bowels, bladder, and digestion. Very small bladders are indicative of the disease, and it is often difficult to diagnose. Death can occur quite suddenly, and some live for only 3–4 years; the oldest recorded was a female cat named Pharrah at 7 years when affected with the disease. In one report, it was shown to affect about 30% of Manx cats studied, but nearly all of those cases were rumpies, which exhibit the most extreme phenotype. Feline expert Roger Tabor has stated: "Only the fact that the Manx is a historic breed stops us being as critical of this dangerous gene as of other more recent selected abnormalities."
The breed is also predisposed to rump fold intertrigo, and to corneal dystrophy. In a review of over 5,000 cases of urate urolithiasis the Manx was noticeably under-represented, with an odds ratio of 0.35.
Some tailless cats such as the Manx cats may develop megacolon, which is a recurring condition causing constipation that can be life-threatening to the cat if not properly monitored. It is a condition in which, due to absence of a tail, the smooth muscle that normally contracts to push stools toward the rectum loses its ability to do so.
Following on updated genetic research, both the Australian Cat Federation and (less stringently) the GCCF impose special breeding restrictions on Manx cats (and derived stock like the Cymric), for animal welfare reasons.
Identification of the Manx Cat tailless gene
In 2013, prior to initiation of the Manx Cat Genome Project (below), genetic mutations in the brachyury gene were shown to be responsible for failure of tail development in the Manx cat, as well as four other tailless breeds of cat. Mutations in orthologs of this gene have been shown to cause tail-loss defects in a number of other species, notably the mouse. Mutations in the human version of the brachyury gene are associated with a range of neural tube defects.
Manx Cat Genome Project
To better understand the genetics of the breed, the Manx Cat Genome Project (MCGP) was launched in August 2015, as a crowdfunded volunteer project by computational biologist Rachel Glover of Douglas, Isle of Man, to perform the first whole genome sequencing of the Manx cat, uncovering the genetic mutations that make the Manx distinct from other cat populations, and to contribute data to the genome databases at the 99 Lives Cat Genome Sequencing Project of the University of Missouri, and the US National Center for Biotechnology Information (NCBI). It is the Isle of Man's first gene sequencing programme, with samples collected and data analysed by MCGP in the Isle of Man, with the input of scientists around the world, initial sequencing work being performed by the firm Edinburgh Genomics and the University of Edinburgh in Scotland and by 99 Lives, and server resources donated by Isle of Man biomedical information technology company ServiceTech.
The project aims to answer four questions:
Which mutations are unique to the breed, aside from the obvious suppressed tail?
What genes are involved in Manx syndrome?
What genes control tail length? (The Manx taillessness gene only determines whether the tail will be suppressed, not the extent of suppression.)
Is there a genetic basis for any health problems associated with the breed other than Manx syndrome?
One desired result of this research is the development of tests that can be used to keep the breed healthy by identifying cats which should not be bred. A minimum of three cats' genes will have to be sequenced to obtain the required genetic data.
After the initial fundraising goal was reached in December 2015, the first cat sequenced was a purebred Manx calico rumpy named Bonnag, selected because the registry of this dam (breeding female) and her kittens in the British Governing Council of the Cat Fancy (GCCF) aids controlled study of a specific bloodline. Bonnag's samples were sent for sequencing in April 2016, with raw gene sequence results received by MCGP in August 2016; the laborious process of genome assembly has begun, to be followed by comparison with previously collected cat genomic data from 99 Lives, and eventual peer-reviewed publication of the results in a scientific journal. Fundraising for the second genome to be sequenced by the project began September 2016; costs dropped to UK£1,400 per cat in November 2015, and as of April 2016 dropped to about £1,200, using the Illumina HiSeq X Ten sequencer, down from original projections of £10,000 before the X Ten was available for non-human sequencing. The dramatic drop in costs allowed the first cat's sequencing to be done well ahead of the original schedule. MCGP has already identified the location of the mutation responsible for suppression of Bonnag's tail, the deletion of a single bit of genetic data among 2.8 billion making up the genome.
The selected second sample is from a kitten that had to be euthanised for Manx syndrome, and it is hoped that this new sequence can identify the genetic specifics of the condition and why it only affects some offspring.
Behaviour
Fancier's often describe the Manx as being doglike in behaviour.
These beliefs about the Manx's behaviour were not described in the past.
Lane's early and experienced account of the temperament of this "variety, which is quaint and interesting" is simply that they were "docile, good-tempered and sociable", and that a prize specimen should be "an alert, active animal of much power and energetic character."
Manx are prized as hunters, known to take down larger prey (e.g. adult rats) even when they are young, and were thus long in demand for working roles like farm cat (Manx: or 'mouser', from 'mouse') and ship's cat ( or loosely 'scratcher, scratchy-one', from or 'scratching, scratchy, scraping').
In popular culture
Isle of Man national symbol
The Isle of Man uses the Manx cat as one of the symbols of the island nation and its unique culture. On Isle of Man currency, Manx cats are the subject of the reverse of four special commemorative crown coins. The first two, issued in 1970 and 1975, are stand-alone releases in both copper-nickel and silver proofs, while the third, in 1988, inaugurated an ongoing series of annual cat coin issues that have also been produced in gold in various sizes; an almost-hidden Manx cat appears in the background on each of the 1989-onward releases featuring other breeds. A Manx, with a kitten, was the featured cat again in 2012.
A Manx cat, in stylized Celtic knotwork art, also appears on the island's 1980–83 penny. The breed figures on numerous Isle of Man postage stamps, including a 2011 series of six that reproduce the art from Victorian era Manx cat postcards, a 1996 one-stamp decorative sheetlet, one stamp in a 1994 tourism 10-stamp booklet, a 1996 five-stamp series of Manx cats around the world, and a 1989 set of the breed in various coat patterns, plus two high-value definitives of 1983 and 1989. The cat appears prominently as the subject of a large number of tourist goods and Manx pride items available on the island and over the Internet, serving (along with the triskelion and the four-horned Manx Loaghtan sheep) as an emblem of the Isle of Man.
Famous real-world Manx cats
All Ball, Lipstick, and Smokey, three Manx cats that were companion animals to Koko, a captive gorilla renowned for communicating in rudimentary American Sign Language
Bob, the male subject of Bob the Preschool Cat: A Biography of an Urban Manx Cat by E. Romayne Hertweck (2009, )
Bonnag, a female Manx, the first of her breed to have her whole genome sequenced (in 2016, by the Manx Cat Genome Project, ), and only the second cat of any breed to receive this level of study (the first was an Abyssinian sequenced by the 99 Lives project in 2014). Bonnag was bred by Zoe Grundey at the Triskele Manx Cats cattery in Douglas, Isle of Man.
Peta, Chief Mouser to the Cabinet Office of the United Kingdom government between 1964 and sometime between 1969 and 1976.
Fictional Manx cats
Bluebeard, from the German animated film Felidae (1994)
Gordon from the American animated TV series Catscratch (2005–2007)
Ma Manx, matriarch of a crime gang in the children's novel Rex Tabby: Cat Detective by Daniel Kirk (2004, )
Mac Manc McManx, a recurring guest character in the American daily comic strip Get Fuzzy (1999–present)
Mayor Manx from the American animated TV series SWAT Kats (1993–1995)
Manx, the antagonist to Slimer, of Slimer! and the Real Ghostbusters (1988–1991)
Manx Cat, the antagonist for the bulk of Paul Gallico's children's novel Manxmouse: The Mouse Who Knew No Fear (1968, ), and the 1979 Japanese anime based on it
Marco the Manx from Joann Roe's series of children's books, Fisherman Cat (1988, ), Castaway Cat (1989, ), Alaska Cat (1990, ), and Samurai Cat (1993, )
Max from Adam Whitmore's "Max the Cat" 1986 series of children's books, Max Leaves Home (), Max in America (), Max in India (), and Max in Australia ()
Mika, title character of the children's book Mika the Manx Cat by D. M. Hart (2012, )
Narrator, an orange Manx, in the children's book The Cats of Grand Central by Laura Archibald, illustrated by Garner Beckett (2003, )
Olaf, protagonist of Olaf Comes Home by Kathy Dollina Creamer (2001, ), a children's book modeled on "The Ugly Duckling"
Raffles, Bernie Rhodenbarr's Manx cat in Lawrence Block's "Burglar" series of mystery novels, first appearing in The Burglar Who Traded Ted Williams. Bernie is not convinced the cat is a Manx but it does have no tail. (1994, )
Stimpy, one of the two main characters of the American animated TV series The Ren and Stimpy Show (1991–1996)
Tiara Boobowski was planned to be a Manx cat character in the Sonic the Hedgehog game Sonic X-Treme but the game was cancelled.
The Manx cat that the narrator sees during a lunch party in chapter one of Virginia Woolf's A Room of One's Own (1929)
Other
The Norton Manx motorcycle line (1947–1962, Norton Motors Ltd.), though ostensibly named after the Isle of Man TT road race (which the brand dominated for decades, until the 1970s), was long promoted with Manx cat badges, in the forms of both enameled metal pins and sew-on patches. The Manx Norton has experienced a major revival among modern enthusiasts of classic motorcycle racing.
The Meyers Manx (1964–1971, B. F. Meyers & Co.) is the original, much-copied Volkswagen Beetle–based dune buggy, and broke desert racing records shortly after its introduction. It was named after the cat, due to its design – short-bodied, tall-wheeled, and manoeuvrable. The original designer has revived and updated it as the "Manxter" (2000–present, Meyers Manx, Inc.).
A popular flying model aircraft of the late 1950s was the Manx Cat, sold in kit form as the Manx Cat V, and in printed plan form as the Manx Cat I through IV, with progressively larger wings. Designed by Bob Buragas, the hand-launched biplane model is constructed of balsa wood, features a very short tail (thus the name), has a 32.5 inch wingspan (in versions IV and V), can accommodate .19 to .35 engine sizes, and can be modified with a Dumas Spectrum "combat" wing. It was profiled in hobbyist magazines, like the February 1957 Flying Models (which details the history of the different models, including a miniature Manx Kitten version), and the October 1958 American Modeler.
A Grimjack comic book story, The Manx Cat, was serialised as a Comicmix.com webcomic in January 2011, and has since seen print as a six-issue miniseries by IDW Comics. The story involves "The Manx Cat", a statuette of such a cat that at first seems to be a simple MacGuffin like the classic Maltese Falcon of the novel and films of that name, but which begins showing malevolent powers. The plot thickens with time travel, reincarnation, and Cthulhu Mythos-style "elder gods". Like most modern comics, it features digitally-colored art, over hand-drawn pencil work.
In popular music, Florrie Forde released a 1930 recording of a Dan Leno Jr comedic music hall song, "What Happened to the Manx Cat's Tail?", as the B-side of "Stein! Stein! Ev'rywhere We Go", on an 8-inch, 78 RPM gramophone record (serial number 1430 on the Edison Bell Radio label).
| Biology and health sciences | Cats | Animals |
145694 | https://en.wikipedia.org/wiki/Oceanic%20trench | Oceanic trench | Oceanic trenches are prominent, long, narrow topographic depressions of the ocean floor. They are typically wide and below the level of the surrounding oceanic floor, but can be thousands of kilometers in length. There are about of oceanic trenches worldwide, mostly around the Pacific Ocean, but also in the eastern Indian Ocean and a few other locations. The greatest ocean depth measured is in the Challenger Deep of the Mariana Trench, at a depth of below sea level.
Oceanic trenches are a feature of the Earth's distinctive plate tectonics. They mark the locations of convergent plate boundaries, along which lithospheric plates move towards each other at rates that vary from a few millimeters to over ten centimeters per year. Oceanic lithosphere moves into trenches at a global rate of about per year. A trench marks the position at which the flexed, subducting slab begins to descend beneath another lithospheric slab. Trenches are generally parallel to and about from a volcanic arc.
Much of the fluid trapped in sediments of the subducting slab returns to the surface at the oceanic trench, producing mud volcanoes and cold seeps. These support unique biomes based on chemotrophic microorganisms. There is concern that plastic debris is accumulating in trenches and threatening these communities.
Geographic distribution
There are approximately of convergent plate margins worldwide. These are mostly located around the Pacific Ocean, but are also found in the eastern Indian Ocean, with a few shorter convergent margin segments in other parts of the Indian Ocean, in the Atlantic Ocean, and in the Mediterranean. They are found on the oceanward side of island arcs and Andean-type orogens. Globally, there are over 50 major ocean trenches covering an area of 1.9 million km2 or about 0.5% of the oceans.
Trenches are geomorphologically distinct from troughs. Troughs are elongated depressions of the sea floor with steep sides and flat bottoms, while trenches are characterized by a V-shaped profile. Trenches that are partially infilled are sometimes described as troughs, for example the Makran Trough. Some trenches are completely buried and lack bathymetric expression as in the Cascadia subduction zone, which is completely filled with sediments. Despite their appearance, in these instances the fundamental plate-tectonic structure is still an oceanic trench. Some troughs look similar to oceanic trenches but possess other tectonic structures. One example is the Lesser Antilles Trough, which is the forearc basin of the Lesser Antilles subduction zone. Also not a trench is the New Caledonia trough, which is an extensional sedimentary basin related to the Tonga-Kermadec subduction zone. Additionally, the Cayman Trough, which is a pull-apart basin within a transform fault zone, is not an oceanic trench.
Trenches, along with volcanic arcs and Wadati–Benioff zones (zones of earthquakes under a volcanic arc) are diagnostic of convergent plate boundaries and their deeper manifestations, subduction zones. Here, two tectonic plates are drifting into each other at a rate of a few millimeters to over per year. At least one of the plates is oceanic lithosphere, which plunges under the other plate to be recycled in the Earth's mantle.
Trenches are related to, but distinct from, continental collision zones, such as the Himalayas. Unlike in trenches, in continental collision zones continental crust enters a subduction zone. When buoyant continental crust enters a trench, subduction comes to a halt and the area becomes a zone of continental collision. Features analogous to trenches are associated with collision zones. One such feature is the peripheral foreland basin, a sediment-filled foredeep. Examples of peripheral foreland basins include the floodplains of the Ganges River and the Tigris-Euphrates river system.
History of the term "trench"
Trenches were not clearly defined until the late 1940s and 1950s. The bathymetry of the ocean was poorly known prior to the Challenger expedition of 1872–1876, which took 492 soundings of the deep ocean. At station #225, the expedition discovered Challenger Deep, now known to be the southern end of the Mariana Trench. The laying of transatlantic telegraph cables on the seafloor between the continents during the late 19th and early 20th centuries provided further motivation for improved bathymetry. The term trench, in its modern sense of a prominent elongated depression of the sea bottom, was first used by Johnstone in his 1923 textbook An Introduction to Oceanography.
During the 1920s and 1930s, Felix Andries Vening Meinesz measured gravity over trenches using a newly developed gravimeter that could measure gravity from aboard a submarine. He proposed the tectogene hypothesis to explain the belts of negative gravity anomalies that were found near island arcs. According to this hypothesis, the belts were zones of downwelling of light crustal rock arising from subcrustal convection currents. The tectogene hypothesis was further developed by Griggs in 1939, using an analogue model based on a pair of rotating drums. Harry Hammond Hess substantially revised the theory based on his geological analysis.
World War II in the Pacific led to great improvements of bathymetry, particularly in the western Pacific. In light of these new measurements, the linear nature of the deeps became clear. There was a rapid growth of deep sea research efforts, especially the widespread use of echosounders in the 1950s and 1960s. These efforts confirmed the morphological utility of the term "trench." Important trenches were identified, sampled, and mapped via sonar.
The early phase of trench exploration reached its peak with the 1960 descent of the Bathyscaphe Trieste to the bottom of the Challenger Deep. Following Robert S. Dietz' and Harry Hess' promulgation of the seafloor spreading hypothesis in the early 1960s and the plate tectonic revolution in the late 1960s, the oceanic trench became an important concept in plate tectonic theory.
Morphology
Oceanic trenches are wide and have an asymmetric V-shape, with the steeper slope (8 to 20 degrees) on the inner (overriding) side of the trench and the gentler slope (around 5 degrees) on the outer (subducting) side of the trench. The bottom of the trench marks the boundary between the subducting and overriding plates, known as the basal plate boundary shear or the subduction décollement. The depth of the trench depends on the starting depth of the oceanic lithosphere as it begins its plunge into the trench, the angle at which the slab plunges, and the amount of sedimentation in the trench. Both starting depth and subduction angle are greater for older oceanic lithosphere, which is reflected in the deep trenches of the western Pacific. Here the bottoms of the Marianas and the Tonga–Kermadec trenches are up to below sea level. In the eastern Pacific, where the subducting oceanic lithosphere is much younger, the depth of the Peru-Chile trench is around .
Though narrow, oceanic trenches are remarkably long and continuous, forming the largest linear depressions on earth. An individual trench can be thousands of kilometers long. Most trenches are convex towards the subducting slab, which is attributed to the spherical geometry of the Earth.
The trench asymmetry reflects the different physical mechanisms that determine the inner and outer slope angle. The outer slope angle of the trench is determined by the bending radius of the subducting slab, as determined by its elastic thickness. Since oceanic lithosphere thickens with age, the outer slope angle is ultimately determined by the age of the subducting slab. The inner slope angle is determined by the angle of repose of the overriding plate edge. This reflects frequent earthquakes along the trench that prevent oversteepening of the inner slope.
As the subducting plate approaches the trench, it bends slightly upwards before beginning its plunge into the depths. As a result, the outer trench slope is bounded by an outer trench high. This is subtle, often only tens of meters high, and is typically located a few tens of kilometers from the trench axis. On the outer slope itself, where the plate begins to bend downwards into the trench, the upper part of the subducting slab is broken by bending faults that give the outer trench slope a horst and graben topography. The formation of these bending faults is suppressed where oceanic ridges or large seamounts are subducting into the trench, but the bending faults cut right across smaller seamounts. Where the subducting slab is only thinly veneered with sediments, the outer slope will often show seafloor spreading ridges oblique to the horst and graben ridges.
Sedimentation
Trench morphology is strongly modified by the amount of sedimentation in the trench. This varies from practically no sedimentation, as in the Tonga-Kermadec trench, to completely filled with sediments, as with the Cascadia subduction zone. Sedimentation is largely controlled by whether the trench is near a continental sediment source. The range of sedimentation is well illustrated by the Chilean trench. The north Chile portion of the trench, which lies along the Atacama Desert with its very slow rate of weathering, is sediment-starved, with from 20 to a few hundred meters of sediments on the trench floor. The tectonic morphology of this trench segment is fully exposed on the ocean bottom. The central Chile segment of the trench is moderately sedimented, with sediments onlapping onto pelagic sediments or ocean basement of the subducting slab, but the trench morphology is still clearly discernible. The southern Chile segment of the trench is fully sedimented, to the point where the outer rise and slope are no longer discernible. Other fully sedimented trenches include the Makran Trough, where sediments are up to thick; the Cascadia subduction zone, which is completed buried by of sediments; and the northernmost Sumatra subduction zone, which is buried under of sediments.
Sediments are sometimes transported along the axis of an oceanic trench. The central Chile trench experiences transport of sediments from source fans along an axial channel. Similar transport of sediments has been documented in the Aleutian trench.
In addition to sedimentation from rivers draining into a trench, sedimentation also takes place from landslides on the tectonically steepened inner slope, often driven by megathrust earthquakes. The Reloca Slide of the central Chile trench is an example of this process.
Erosive versus accretionary margins
Convergent margins are classified as erosive or accretionary, and this has a strong influence on the morphology of the inner slope of the trench. Erosive margins, such as the northern Peru-Chile, Tonga-Kermadec, and Mariana trenches, correspond to sediment-starved trenches. The subducting slab erodes material from the lower part of the overriding slab, reducing its volume. The edge of the slab experiences subsidence and steepening, with normal faulting. The slope is underlain by relative strong igneous and metamorphic rock, which maintains a high angle of repose. Over half of all convergent margins are erosive margins.
Accretionary margins, such as the southern Peru-Chile, Cascadia, and Aleutians, are associated with moderately to heavily sedimented trenches. As the slab subducts, sediments are "bulldozed" onto the edge of the overriding plate, producing an accretionary wedge or accretionary prism. This builds the overriding plate outwards. Because the sediments lack strength, their angle of repose is gentler than the rock making up the inner slope of erosive margin trenches. The inner slope is underlain by imbricated thrust sheets of sediments. The inner slope topography is roughened by localized mass wasting. Cascadia has practically no bathymetric expression of the outer rise and trench, due to complete sediment filling, but the inner trench slope is complex, with many thrust ridges. These compete with canyon formation by rivers draining into the trench. Inner trench slopes of erosive margins rarely show thrust ridges.
Accretionary prisms grow in two ways. The first is by frontal accretion, in which sediments are scraped off the downgoing plate and emplaced at the front of the accretionary prism. As the accretionary wedge grows, older sediments further from the trench become increasingly lithified, and faults and other structural features are steepened by rotation towards the trench. The other mechanism for accretionary prism growth is underplating (also known as basal accretion) of subducted sediments, together with some oceanic crust, along the shallow parts of the subduction decollement. The Franciscan Group of California is interpreted as an ancient accretionary prism in which underplating is recorded as tectonic mélanges and duplex structures.
Earthquakes
Frequent megathrust earthquakes modify the inner slope of the trench by triggering massive landslides. These leave semicircular landslide scarps with slopes of up to 20 degrees on the headwalls and sidewalls.
Subduction of seamounts and aseismic ridges into the trench may increase aseismic creep and reduce the severity of earthquakes. Contrariwise, subduction of large amounts of sediments may allow ruptures along the subduction décollement to propagate for great distances to produce megathrust earthquakes.
Trench rollback
Trenches seem positionally stable over time, but scientists believe that some trenches—particularly those associated with subduction zones where two oceanic plates converge—move backward into the subducting plate. This is called trench rollback or hinge retreat (also hinge rollback) and is one explanation for the existence of back-arc basins.
Forces perpendicular to the slab (the portion of the subducting plate within the mantle) are responsible for steepening of the slab and, ultimately, the movement of the hinge and trench at the surface. These forces arise from the negative buoyancy of the slab with respect to the mantle modified by the geometry of the slab itself. The extension in the overriding plate, in response to the subsequent subhorizontal mantle flow from the displacement of the slab, can result in formation of a back-arc basin.
Processes involved
Several forces are involved in the process of slab rollback. Two forces acting against each other at the interface of the two subducting plates exert forces against one another. The subducting plate exerts a bending force (FPB) that supplies pressure during subduction, while the overriding plate exerts a force against the subducting plate (FTS). The slab pull force (FSP) is caused by the negative buoyancy of the plate driving the plate to greater depths. The resisting force from the surrounding mantle opposes the slab pull forces. Interactions with the 660-km discontinuity cause a deflection due to the buoyancy at the phase transition (F660). The unique interplay of these forces is what generates slab rollback. When the deep slab section obstructs the down-going motion of the shallow slab section, slab rollback occurs. The subducting slab undergoes backward sinking due to the negative buoyancy forces causing a retrogradation of the trench hinge along the surface. Upwelling of the mantle around the slab can create favorable conditions for the formation of a back-arc basin.
Seismic tomography provides evidence for slab rollback. Results demonstrate high temperature anomalies within the mantle suggesting subducted material is present in the mantle. Ophiolites are viewed as evidence for such mechanisms as high pressure and temperature rocks are rapidly brought to the surface through the processes of slab rollback, which provides space for the exhumation of ophiolites.
Slab rollback is not always a continuous process suggesting an episodic nature. The episodic nature of the rollback is explained by a change in the density of the subducting plate, such as the arrival of buoyant lithosphere (a continent, arc, ridge, or plateau), a change in the subduction dynamics, or a change in the plate kinematics. The age of the subducting plates does not have any effect on slab rollback. Nearby continental collisions have an effect on slab rollback. Continental collisions induce mantle flow and extrusion of mantle material, which causes stretching and arc-trench rollback. In the area of the Southeast Pacific, there have been several rollback events resulting in the formation of numerous back-arc basins.
Mantle interactions
Interactions with the mantle discontinuities play a significant role in slab rollback. Stagnation at the 660-km discontinuity causes retrograde slab motion due to the suction forces acting at the surface. Slab rollback induces mantle return flow, which causes extension from the shear stresses at the base of the overriding plate. As slab rollback velocities increase, circular mantle flow velocities also increase, accelerating extension rates. Extension rates are altered when the slab interacts with the discontinuities within the mantle at 410 km and 660 km depth. Slabs can either penetrate directly into the lower mantle, or can be retarded due to the phase transition at 660 km depth creating a difference in buoyancy. An increase in retrograde trench migration (slab rollback) (2–4 cm/yr) is a result of flattened slabs at the 660-km discontinuity where the slab does not penetrate into the lower mantle. This is the case for the Japan, Java and Izu–Bonin trenches. These flattened slabs are only temporarily arrested in the transition zone. The subsequent displacement into the lower mantle is caused by slab pull forces, or the destabilization of the slab from warming and broadening due to thermal diffusion. Slabs that penetrate directly into the lower mantle result in slower slab rollback rates (~1–3 cm/yr) such as the Mariana arc, Tonga arcs.
Hydrothermal activity and associated biomes
As sediments are subducted at the bottom of trenches, much of their fluid content is expelled and moves back along the subduction décollement to emerge on the inner slope as mud volcanoes and cold seeps. Methane clathrates and gas hydrates also accumulate in the inner slope, and there is concern that their breakdown could contribute to global warming.
The fluids released at mud volcanoes and cold seeps are rich in methane and hydrogen sulfide, providing chemical energy for chemotrophic microorganisms that form the base of a unique trench biome. Cold seep communities have been identified in the inner trench slopes of the western Pacific (especially Japan), South America, Barbados, the Mediterranean, Makran, and the Sunda trench. These are found at depths as great as . The genome of the extremophile Deinococcus from Challenger Deep has sequenced for its ecological insights and potential industrial uses.
Because trenches are the lowest points in the ocean floor, there is concern that plastic debris may accumulate in trenches and endanger the fragile trench biomes.
Deepest oceanic trenches
Recent measurements, where the salinity and temperature of the water was measured throughout the dive, have uncertainties of about . Older measurements may be off by hundreds of meters.
Notable oceanic trenches
(*) The five deepest trenches in the world
Ancient oceanic trenches
| Physical sciences | Oceanic and coastal landforms | null |
145700 | https://en.wikipedia.org/wiki/Crust%20%28geology%29 | Crust (geology) | In geology, the crust is the outermost solid shell of a planet, dwarf planet, or natural satellite. It is usually distinguished from the underlying mantle by its chemical makeup; however, in the case of icy satellites, it may be defined based on its phase (solid crust vs. liquid mantle).
The crusts of Earth, Mercury, Venus, Mars, Io, the Moon and other planetary bodies formed via igneous processes and were later modified by erosion, impact cratering, volcanism, and sedimentation.
Most terrestrial planets have fairly uniform crusts. Earth, however, has two distinct types: continental crust and oceanic crust. These two types have different chemical compositions and physical properties and were formed by different geological processes.
Types of crust
Planetary geologists divide crust into three categories based on how and when it formed.
Primary crust / primordial crust
This is a planet's "original" crust. It forms from solidification of a magma ocean. Toward the end of planetary accretion, the terrestrial planets likely had surfaces that were magma oceans. As these cooled, they solidified into crust. This crust was likely destroyed by large impacts and re-formed many times as the Era of Heavy Bombardment drew to a close.
The nature of primary crust is still debated: its chemical, mineralogic, and physical properties are unknown, as are the igneous mechanisms that formed them. This is because it is difficult to study: none of Earth's primary crust has survived to today. Earth's high rates of erosion and crustal recycling from plate tectonics has destroyed all rocks older than about 4 billion years, including whatever primary crust Earth once had.
However, geologists can glean information about primary crust by studying it on other terrestrial planets. Mercury's highlands might represent primary crust, though this is debated. The anorthosite highlands of the Moon are primary crust, formed as plagioclase crystallized out of the Moon's initial magma ocean and floated to the top; however, it is unlikely that Earth followed a similar pattern, as the Moon was a water-less system and Earth had water. The Martian meteorite ALH84001 might represent primary crust of Mars; however, again, this is debated. Like Earth, Venus lacks primary crust, as the entire planet has been repeatedly resurfaced and modified.
Secondary crust
Secondary crust is formed by partial melting of mostly silicate materials in the mantle, and so is usually basaltic in composition.
This is the most common type of crust in the Solar System. Most of the surfaces of Mercury, Venus, Earth, and Mars comprise secondary crust, as do the lunar maria. On Earth secondary crust forms primarily at mid-ocean spreading centers, where the adiabatic rise of mantle causes partial melting.
Tertiary crust
Tertiary crust is more chemically-modified than either primary or secondary. It can form in several ways:
Igneous processes: partial-melting of secondary crust, coupled with differentiation or dehydration
Erosion and sedimentation: sediments derived from primary, secondary, or tertiary crust
The only known example of tertiary crust is the continental crust of the Earth. It is unknown whether other terrestrial planets can be said to have tertiary crust, though the evidence so far suggests that they do not. This is likely because plate tectonics is needed to create tertiary crust, and Earth is the only planet in the Solar System with plate tectonics.
Earth's crust
Earth's crust is a thin shell on the outside of Earth, accounting for less than 1% of Earth's volume. It is the top component of the lithosphere, a division of Earth's layers that includes the crust and the upper part of the mantle. The lithosphere is broken into tectonic plates that move, allowing heat to escape from the interior of Earth into space.
Moon's crust
A theoretical protoplanet named "Theia" is thought to have collided with the forming Earth, and part of the material ejected into space by the collision accreted to form the Moon. As the Moon formed, the outer part of it is thought to have been molten, a "lunar magma ocean". Plagioclase feldspar crystallized in large amounts from this magma ocean and floated toward the surface. The cumulate rocks form much of the crust. The upper part of the crust probably averages about 88% plagioclase (near the lower limit of 90% defined for anorthosite): the lower part of the crust may contain a higher percentage of ferromagnesian minerals such as the pyroxenes and olivine, but even that lower part probably averages about 78% plagioclase. The underlying mantle is denser and olivine-rich.
The thickness of the crust ranges between about 20 and 120 km. Crust on the far side of the Moon averages about 12 km thicker than that on the near side. Estimates of average thickness fall in the range from about 50 to 60 km. Most of this plagioclase-rich crust formed shortly after formation of the Moon, between about 4.5 and 4.3 billion years ago. Perhaps 10% or less of the crust consists of igneous rock added after the formation of the initial plagioclase-rich material. The best-characterized and most voluminous of these later additions are the mare basalts formed between about 3.9 and 3.2 billion years ago. Minor volcanism continued after 3.2 billion years, perhaps as recently as 1 billion years ago. There is no evidence of plate tectonics.
Study of the Moon has established that a crust can form on a rocky planetary body significantly smaller than Earth. Although the radius of the Moon is only about a quarter that of Earth, the lunar crust has a significantly greater average thickness. This thick crust formed almost immediately after formation of the Moon. Magmatism continued after the period of intense meteorite impacts ended about 3.9 billion years ago, but igneous rocks younger than 3.9 billion years make up only a minor part of the crust.
| Physical sciences | Geology | null |
145703 | https://en.wikipedia.org/wiki/Laurasia | Laurasia | Laurasia () was the more northern of two large landmasses that formed part of the Pangaea supercontinent from around (Mya), the other being Gondwana. It separated from Gondwana (beginning in the late Triassic period) during the breakup of Pangaea, drifting farther north after the split and finally broke apart with the opening of the North Atlantic Ocean c. 56 Mya. The name is a portmanteau of Laurentia and Eurasia.
Laurentia, Avalonia, Baltica, and a series of smaller terranes, collided in the Caledonian orogeny c. 400 Ma to form Laurussia. Laurussia then collided with Gondwana to form Pangaea. Kazakhstania and Siberia were then added to Pangaea 290–300 Ma to form Laurasia. Laurasia finally became an independent continental mass when Pangaea broke up into Gondwana and Laurasia.
Terminology and origin of the concept
Laurentia, the Palaeozoic core of North America and continental fragments that now make up part of Europe, collided with Baltica and Avalonia in the Caledonian orogeny from c. 430–420 Mya to form Laurussia. In the Late Carboniferous, Laurussia and Gondwana collided and formed Pangaea. Siberia and Kazakhstania finally collided with Baltica in the Late Permian to form Laurasia. A series of continental blocks that now form East Asia and Southeast Asia were later added to Laurasia.
In 1904–1909, Austrian geologist Eduard Suess proposed that the continents in the Southern Hemisphere were once merged into a larger continent called Gondwana. In 1915, German meteorologist Alfred Wegener proposed the existence of a supercontinent that he called Pangaea. In 1937, South African geologist Alexander du Toit proposed that Pangaea was divided into two larger landmasses, Laurasia in the Northern Hemisphere and Gondwana in the Southern Hemisphere, separated by the Tethys Ocean.
"Laurussia" was defined by Swiss geologist Peter Ziegler in 1988 as the merger between Laurentia and Baltica along the northern Caledonian suture. The "Old Red Continent" is an informal name often used for the Silurian-Carboniferous deposits in the central landmass of Laurussia.
Several earlier supercontinents proposed and debated in the 1990s and later (e.g. Rodinia, Nuna, Nena) included earlier connections between Laurentia, Baltica, Siberia. These original connections apparently survived through one and possibly even two Wilson Cycles, though their intermittent duration and recurrent fit is debated.
Proto-Laurasia
Pre–Rodinia
Laurentia and Baltica first formed a continental mass known as Proto-Laurasia as part of the supercontinent Columbia which was assembled 2,100—1,800 Mya to encompass virtually all known Archaean continental blocks. Surviving sutures from this assembly are the Trans-Hudson orogen in Laurentia; Nagssugtoqidian orogen in Greenland; the Kola-Karelian (the northwest margin of the Svecofennian orogen) and the Volhyn—Central Russia orogen and Pachelma orogen (across western Russia) in Baltica, the Akitkan Orogen in Siberia.
Additional Proterozoic crust was accreted 1,800—1,300 Mya, especially along the Laurentia—Greenland—Baltica margin. Laurentia and Baltica formed a coherent continental mass with southern Greenland and Labrador adjacent to the Arctic margin of Baltica. A magmatic arc extended from Laurentia through southern Greenland to northern Baltica. The breakup of Columbia began 1,600 Mya, including along the western margin of Laurentia and northern margin of Baltica (modern coordinates), and was completed c. 1,300—1,200 Mya, a period during which mafic dike swarms were emplaced, including MacKenzie and Sudbury in Laurentia.
Traces left by large igneous provinces provide evidences for continental mergers during this period. Those related to Proto-Laurasia includes:
1,750 Mya extensive magmatism in Baltica, Sarmatia (Ukraine), southern Siberia, northern Laurentia, and West Africa indicate these cratons were linked to each other;
a 1,630–1,640 Mya-old continent composed of Siberia, Laurentia, and Baltica is suggested by sills in southern Siberia that can be connected to the Melville Bugt dyke swarm in western Greenland;
a major large igneous province 1,380 Mya during the breakup of the Nuna supercontinent connects Laurentia, Baltica, Siberia, Congo, West Africa.
Rodinia
In the vast majority of plate tectonic reconstructions, Laurentia formed the core of the supercontinent Rodinia, but the exact fit of various continents within Rodinia is debated. In some reconstructions, Baltica was attached to Greenland along its Scandinavian margin while Amazonia was docked along Baltica's Tornquist margin. Australia and East Antarctica were located on Laurentia's western margin.
Siberia was located near but at some distance from Laurentia's northern margin in most reconstructions. In the reconstruction of some Russian geologists, however, the southern margin (modern coordinates) of Siberia merged with the northern margin of Laurentia, and these two continents broke up along what is now the -long Central Asian Foldbelt no later than 570 Mya and traces of this breakup can still be found in the Franklin dike swarm in northern Canada and the Aldan Shield in Siberia.
The Proto-Pacific opened and Rodinia began to breakup during the Neoproterozoic (c. 750–600 Mya) as Australia-Antarctica (East Gondwana) rifted from the western margin of Laurentia, while the rest of Rodinia (West Gondwana and Laurasia) rotated clockwise and drifted south. Earth subsequently underwent a series of glaciations – the Varanger (c. 650 Mya, also known as Snowball Earth) and the Rapitan and Ice Brook glaciations (c. 610-590 Mya) – both Laurentia and Baltica were located south of 30°S, with the South Pole located in eastern Baltica, and glacial deposits from this period have been found in Laurentia and Baltica but not in Siberia.
A mantle plume (the Central Iapetus Magmatic Province) forced Laurentia and Baltica to separate ca. 650–600 Mya and the Iapetus Ocean opened between them. Laurentia then began to move quickly () north towards the Equator where it got stuck over a cold spot in the Proto-Pacific. Baltica remained near Gondwana in southern latitudes into the Ordovician.
Pannotia
Laurentia, Baltica, and Siberia remained connected to each other within the short-lived, Precambrian-Cambrian supercontinent Pannotia or Greater Gondwana. At this time a series of continental blocks called as Peri-Gondwana, that now form parts of Asia, the Cathaysian terranes, namely Indochina, North China, South China , Cimmerian terranes, Sibumasu, Qiangtang, Lhasa, Afghanistan, Iran, Turkey – were still attached to the Indian–Australian margin of Gondwana. Other blocks that now form part of southwestern Europe and North America from New England to Florida were still attached to the African-South American margin of Gondwana. This northward drift of terranes across the Tethys Ocean also included the Hunic terranes, now spread from Europe to China.
Pannotia broke apart in the late Precambrian into Laurentia, Baltica, Siberia, Gondwana. A series of continental blocks, the Cadomian, Avalonian, Cathaysian, Cimmerian terranes, broke away from Gondwana and began to drift north.
Euramerica/Laurussia
Laurentia remained almost static near the Equator throughout the early Palaeozoic, separated from Baltica by the up to -wide Iapetus Ocean. In the Late Cambrian, the mid-ocean ridge in the Iapetus Ocean subducted beneath Gondwana which resulted in the opening of a series of large back-arc basins. During the Ordovician, these basins evolved into a new ocean, the Rheic Ocean, which separated a series of terranes – Avalonia, Carolinia, and Armorica – from Gondwana.
Avalonia rifted from Gondwana in the Early Ordovician and collided with Baltica near the Ordovician–Silurian boundary (480–420 Mya). Baltica-Avalonia was then rotated and pushed north towards Laurentia. The collision between these continents closed the Iapetus Ocean and formed Laurussia, also known as Euramerica. Another historical term for this continent is the Old Red Continent or Old Red Sandstone Continent, in reference to abundant red beds of the Old Red Sandstone during the Devonian. The continent covered including several large Arctic continental blocks.
With the Caledonian orogeny completed Laurussia was delimited thus:
The eastern margin were the Barents Shelf and Moscow Platform;
the western margin were the western shelves of Laurentia, later affected by the Antler orogeny;
the northern margin was the Innuitian-Lomonosov orogeny which marked the collision between Laurussia and the Arctic Craton;
the southern margin was a Pacific-style active margin where the northward directed subduction of the ocean floor between Gondwana and Laurussia pushed continental fragments towards the latter.
During the Devonian (416-359 Mya) the combined landmass of Baltica and Avalonia rotated around Laurentia, which remained static near the Equator. The Laurentian warm, shallow seas and on shelves a diverse assemblage of benthos evolved, including the largest trilobites exceeding . The Old Red Sandstone Continent stretched across northern Laurentia and into Avalonia and Baltica but for most of the Devonian a narrow seaway formed a barrier where the North Atlantic would later open. Tetrapods evolved from fish in the Late Devonian, with the oldest known fossils from Greenland. Low sea-levels during the Early Devonian produced natural barriers in Laurussia which resulted in provincialism within the benthic fauna. In Laurentia the Transcontinental Arch divided brachiopods into two provinces, with one of them confined to a large embayment west of the Appalachians. By the Middle Devonian, these two provinces had been united into one and the closure of the Rheic Ocean finally united faunas across Laurussia. High plankton productivity from the Devonian-Carboniferous boundary resulted in anoxic events that left black shales in the basins of Laurentia.
Pangaea
The subduction of the Iapetus Ocean resulted in the first contact between Laurussia and Gondwana in the Late Devonian and terminated in full collision or the Variscan orogeny in the early Carboniferous (340 Mya). The Variscan orogeny closed the Rheic Ocean (between Avalonia and Armorica) and the Proto-Tethys Ocean (between Armorica and Gondwana) to form the supercontinent Pangaea. The Variscan orogeny is complex and the exact timing and the order of the collisions between involved microcontinents has been debated for decades.
Pangaea was completely assembled by the Permian except for the Asian blocks. The supercontinent was centred on the Equator during the Triassic and Jurassic, a period that saw the emergence of the Pangaean megamonsoon. Heavy rainfall resulted in high groundwater tables, in turn resulting in peat formation and extensive coal deposits.
During the Cambrian and Early Ordovician, when wide oceans separated all major continents, only pelagic marine organisms, such as plankton, could move freely across the open ocean and therefore the oceanic gaps between continents are easily detected in the fossil records of marine bottom dwellers and non-marine species. By the Late Ordovician, when continents were pushed closer together closing the oceanic gaps, benthos (brachiopods and trilobites) could spread between continents while ostracods and fishes remained isolated. As Laurussia formed during the Devonian and Pangaea formed, fish species in both Laurussia and Gondwana began to migrate between continents and before the end of the Devonian similar species were found on both sides of what remained of the Variscan barrier.
The oldest tree fossils are from the Middle Devonian pteridophyte Gilboa Fossil Forest in central Laurussia (today New York City, United States). In the late Carboniferous, Laurussia was centred on the Equator and covered by tropical rainforests, commonly referred to as the coal forests. By the Permian, the climate had become arid and these Carboniferous rainforests collapsed, lycopsids (giant mosses) were replaced by treeferns. In the dry climate a detritivorous fauna – including ringed worms, molluscs, and some arthropods – evolved and diversified, alongside other arthropods who were herbivorous and carnivorous, and tetrapods – insectivores and piscivores such as amphibians and early amniotes.
Laurasia
During the Carboniferous to Permian periods, Siberia, Kazakhstania, Baltica collided in the Uralian orogeny to form Laurasia.
The Palaezoic-Mesozoic transition was marked by the reorganisation of Earth's tectonic plates which resulted in the assembly of Pangaea, and eventually its break-up. Caused by the detachment of subducted mantle slabs, this reorganisation resulted in rising mantle plumes that produced large igneous provinces when they reached the crust. This tectonic activity also resulted in the Permian–Triassic extinction event. Tentional stresses across Eurasia developed into a large system of rift basins (Urengoy, East Uralian-Turgay, Khudosey) and flood basalts in the West Siberian Basin, the Pechora Basin, South China.
Laurasia and Gondwana were equal in size but had distinct geological histories. Gondwana was assembled before the formation of Pangaea, but the assembly of Laurasia occurred during and after the formation of the supercontinent. These differences resulted in different patterns of basin formation and transport of sediments. East Antarctica was the highest ground within Pangaea and produced sediments that were transported across eastern Gondwana but never reached Laurasia. During the Palaeozoic, c. 30–40% of Laurasia was covered by shallow marine water but only 10–20% of Gondwana was covered by shallow marine water.
Asian blocks
During the assembly of Pangaea, Laurasia grew as continental blocks broke off Gondwana's northern margin; pulled by old closing oceans in front of them and pushed by new opening oceans behind them. During the Neoproterozoic-Early Paleozoic break-up of Rodinia, the opening of the Proto-Tethys Ocean split the Asian blocks – Tarim, Qaidam, Alex, North China, South China – from the northern shores of Gondwana (north of India and Australia in modern coordinates) and the closure of the same ocean reassembled them along the same shores 500–460 Mya resulting in Gondwana at its largest extent.
The break-up of Rodinia also resulted in the opening of the long-lived Paleo-Asian Ocean between Baltica and Siberia in the north and Tarim and North China in the south. The closure of this ocean is preserved in the Central Asian Orogenic Belt, the largest orogen on Earth.
North China, South China, Indochina, Tarim broke off from Gondwana during the Silurian to Devonian periods; as the Paleo-Tethys Ocean opened behind them. Sibumasu and Qiantang and other Cimmerian continental fragments broke off in the Early Permian. Lhasa, Burma, Sikuleh, southwest Sumatra, West Sulawesi, and parts of Borneo, broke off during the Late Triassic-Late Jurassic.
During the Carboniferous and Permian, Baltica first collided with Kazakhstania and Siberia, then North China with Mongolia and Siberia. By the middle Carboniferous, however, South China had already been in contact with North China long enough to allow floral exchange between the two continents. The Cimmerian blocks rifted from Gondwana in the Late Carboniferous.
In the early Permian, the Neo-Tethys Ocean opened behind the Cimmerian terranes (Sibumasu, Qiantang, Lhasa) and, in the late Carboniferous, the Paleo-Tethys Ocean closed in front. The eastern branch of the Paleo-Tethys Ocean, however, remained opened while Siberia was added to Laurussia and Gondwana collided with Laurasia.
When the eastern Palaeo-Tethys closed 250–230 Mya, a series of Asian blocks – Sibumasu, Indochina, South China, Qiantang, Lhasa – formed a separate southern Asian continent. This continent collided 240–220 Mya with a northern continent – North China, Qinling, Qilian, Qaidam, Alex, Tarim – along the Central China orogen to form a combined East Asian continent. The northern margins of the northern continent collided with Baltica and Siberia 310–250 Ma, and thus the formation of the East Asian continent marked Pangaea at its greatest extent. By this time, the rifting of western Pangaea had already begun.
Flora and fauna
Pangaea split in two as the Tethys Seaway opened between Gondwana and Laurasia in the Late Jurassic. The fossil record, however, suggests the intermittent presence of a Trans-Tethys land bridge, though the location and duration of such a land bridge remains enigmatic.
Pine trees evolved in the early Mesozoic c. 250 Mya and the pine genus originated in Laurasia in the Early Cretaceous c. 130 Mya in competition with faster growing flowering plants. Pines adapted to cold and arid climates in environments where the growing season was shorter or wildfire common; this evolution limited pine range to between 31° and 50° north and resulted in a split into two subgenera: Strobus adapted to stressful environments and Pinus to fire-prone landscapes. By the end of the Cretaceous, pines were established across Laurasia, from North America to East Asia.
From the Triassic to the Early Jurassic, before the break-up of Pangaea, archosaurs (crurotarsans, pterosaurs and dinosaurs including birds) had a global distribution, especially crurotarsans, the group ancestral to the crocodilians. This cosmopolitanism ended as Gondwana fragmented and Laurasia was assembled. Pterosaur diversity reach a maximum in the Late Jurassic—Early Cretaceous and plate tectonic didn't affect the distribution of these flying reptiles. Crocodilian ancestors also diversified during the Early Cretaceous but were divided into Laurasian and Gondwanan populations; true crocodilians evolved from the former. The distribution of the three major groups of dinosaurs – the sauropods, theropods, and ornithischians – was similar to that of the crocodilians. East Asia remained isolated with endemic species including psittacosaurs (horned dinosaurs) and Ankylosauridae (club-tailed, armoured dinosaurs).
Meanwhile, mammals slowly settled in Laurasia from Gondwana in the Triassic, the latter of which was the living area of their Permian ancestors. They split in two groups, with one returning to Gondwana (and stayed there after Pangaea split) while the other staying in Laurasia (until further descendants switched to Gondwana starting from the Jurassic).
In the early Eocene, a peak in global warming led to a pan-Arctic fauna with alligators and amphibians present north of the Arctic Circle. In the early Paleogene, landbridges still connected continents, allowing land animals to migrate between them. On the other hand, submerged areas occasionally divided continents: the Turgai Strait separated Europe and Asia from the Middle Jurassic to the Oligocene and as this strait dried out, a massive faunal interchange took place and the resulting extinction event in Europe is known as the Grande Coupure.
The Coraciiformes (an order of birds including kingfishers) evolved in Laurasia. While this group now has a mostly tropical distribution, they originated in the Arctic in the late Eocene c. 35 Mya from where they diversified across Laurasia and farther south across the Equator.
The placental mammal group of Laurasiatheria is named after Laurasia.
Final split
In the Triassic–Early Jurassic (c. 200 Mya), the opening of the Central Atlantic Ocean was preceded by the formation of a series of large rift basins, such as the Newark Basin, between eastern North America, from what is today the Gulf of Mexico to Nova Scotia, and in Africa and Europe, from Morocco to Greenland.
By spreading had begun in the North Atlantic between the Rockall Basin, a continental fragment sitting on top of the Eurasian Plate, and North America. By 56 Mya, Greenland had become an independent plate, separated from North America by the Labrador Sea-Baffin Bay Rift. By 33 Mya, spreading had ceased in the Labrador Sea and relocated to the Mid-Atlantic Ridge. The opening of the North Atlantic Ocean had effectively broken Laurasia in two.
| Physical sciences | Paleogeography | Earth science |
145716 | https://en.wikipedia.org/wiki/Lithosphere | Lithosphere | A lithosphere () is the rigid, outermost rocky shell of a terrestrial planet or natural satellite. On Earth, it is composed of the crust and the lithospheric mantle, the topmost portion of the upper mantle that behaves elastically on time scales of up to thousands of years or more. The crust and upper mantle are distinguished on the basis of chemistry and mineralogy.
Earth's lithosphere
Earth's lithosphere, which constitutes the hard and rigid outer vertical layer of the Earth, includes the crust and the lithospheric mantle (or mantle lithosphere), the uppermost part of the mantle that is not convecting. The lithosphere is underlain by the asthenosphere which is the weaker, hotter, and deeper part of the upper mantle that is able to convect. The lithosphere–asthenosphere boundary is defined by a difference in response to stress. The lithosphere remains rigid for very long periods of geologic time in which it deforms elastically and through brittle failure, while the asthenosphere deforms viscously and accommodates strain through plastic deformation.
The thickness of the lithosphere is thus considered to be the depth to the isotherm associated with the transition between brittle and viscous behavior. The temperature at which olivine becomes ductile (~) is often used to set this isotherm because olivine is generally the weakest mineral in the upper mantle.
The lithosphere is subdivided horizontally into tectonic plates, which often include terranes accreted from other plates.
History of the concept
The concept of the lithosphere as Earth's strong outer layer was described by the English mathematician A. E. H. Love in his 1911 monograph "Some problems of Geodynamics" and further developed by the American geologist Joseph Barrell, who wrote a series of papers about the concept and introduced the term "lithosphere". The concept was based on the presence of significant gravity anomalies over continental crust, from which he inferred that there must exist a strong, solid upper layer (which he called the lithosphere) above a weaker layer which could flow (which he called the asthenosphere). These ideas were expanded by the Canadian geologist Reginald Aldworth Daly in 1940 with his seminal work "Strength and Structure of the Earth." They have been broadly accepted by geologists and geophysicists. These concepts of a strong lithosphere resting on a weak asthenosphere are essential to the theory of plate tectonics.
Types
The lithosphere can be divided into oceanic and continental lithosphere. Oceanic lithosphere is associated with oceanic crust (having a mean density of about ) and exists in the ocean basins. Continental lithosphere is associated with continental crust (having a mean density of about ) and underlies the continents and continental shelves.
Oceanic lithosphere
Oceanic lithosphere consists mainly of mafic crust and ultramafic mantle (peridotite) and is denser than continental lithosphere. Young oceanic lithosphere, found at mid-ocean ridges, is no thicker than the crust, but oceanic lithosphere thickens as it ages and moves away from the mid-ocean ridge. The oldest oceanic lithosphere is typically about thick. This thickening occurs by conductive cooling, which converts hot asthenosphere into lithospheric mantle and causes the oceanic lithosphere to become increasingly thick and dense with age. In fact, oceanic lithosphere is a thermal boundary layer for the convection in the mantle. The thickness of the mantle part of the oceanic lithosphere can be approximated as a thermal boundary layer that thickens as the square root of time.
Here, is the thickness of the oceanic mantle lithosphere, is the thermal diffusivity (approximately ) for silicate rocks, and is the age of the given part of the lithosphere. The age is often equal to L/V, where L is the distance from the spreading centre of mid-ocean ridge, and V is velocity of the lithospheric plate.
Oceanic lithosphere is less dense than asthenosphere for a few tens of millions of years but after this becomes increasingly denser than asthenosphere. While chemically differentiated oceanic crust is lighter than asthenosphere, thermal contraction of the mantle lithosphere makes it more dense than the asthenosphere. The gravitational instability of mature oceanic lithosphere has the effect that at subduction zones, oceanic lithosphere invariably sinks underneath the overriding lithosphere, which can be oceanic or continental. New oceanic lithosphere is constantly being produced at mid-ocean ridges and is recycled back to the mantle at subduction zones. As a result, oceanic lithosphere is much younger than continental lithosphere: the oldest oceanic lithosphere is about 170 million years old, while parts of the continental lithosphere are billions of years old.
Subducted lithosphere
Geophysical studies in the early 21st century posit that large pieces of the lithosphere have been subducted into the mantle as deep as to near the core-mantle boundary, while others "float" in the upper mantle. Yet others stick down into the mantle as far as but remain "attached" to the continental plate above, similar to the extent of the old concept of "tectosphere" revisited by Jordan in 1988. Subducting lithosphere remains rigid (as demonstrated by deep earthquakes along Wadati–Benioff zone) to a depth of about .
Continental lithosphere
Continental lithosphere has a range in thickness from about to perhaps ; the upper approximately of typical continental lithosphere is crust. The crust is distinguished from the upper mantle by the change in chemical composition that takes place at the Moho discontinuity. The oldest parts of continental lithosphere underlie cratons, and the mantle lithosphere there is thicker and less dense than typical; the relatively low density of such mantle "roots of cratons" helps to stabilize these regions.
Because of its relatively low density, continental lithosphere that arrives at a subduction zone cannot subduct much further than about before resurfacing. As a result, continental lithosphere is not recycled at subduction zones the way oceanic lithosphere is recycled. Instead, continental lithosphere is a nearly permanent feature of the Earth.
Mantle xenoliths
Geoscientists can directly study the nature of the subcontinental mantle by examining mantle xenoliths brought up in kimberlite, lamproite, and other volcanic pipes. The histories of these xenoliths have been investigated by many methods, including analyses of abundances of isotopes of osmium and rhenium. Such studies have confirmed that mantle lithospheres below some cratons have persisted for periods in excess of 3 billion years, despite the mantle flow that accompanies plate tectonics.
Microorganisms
The upper part of the lithosphere is a large habitat for microorganisms, with some found more than below Earth's surface.
| Physical sciences | Geology | null |
145717 | https://en.wikipedia.org/wiki/Asthenosphere | Asthenosphere | The asthenosphere () is the mechanically weak and ductile region of the upper mantle of Earth. It lies below the lithosphere, at a depth between c. below the surface, and extends as deep as . However, the lower boundary of the asthenosphere is not well defined.
The asthenosphere is almost solid, but a slight amount of melting (less than 0.1% of the rock) contributes to its mechanical weakness. More extensive decompression melting of the asthenosphere takes place where it wells upwards, and this is the most important source of magma on Earth. It is the source of mid-ocean ridge basalt (MORB) and of some magmas that erupt above subduction zones or in regions of continental rifting.
Characteristics
The asthenosphere is a part of the upper mantle just below the lithosphere that is involved in plate tectonic movement and isostatic adjustments. It is composed of peridotite, a rock containing mostly the minerals olivine and pyroxene. The lithosphere-asthenosphere boundary is conventionally taken at the isotherm. Closer to the surface at lower temperatures, the mantle behaves rigidly; deeper below the surface at higher temperatures, the mantle moves in a ductile fashion. The asthenosphere is where the mantle rock most closely approaches its melting point, and a small amount of melt is likely to present in this layer.
Seismic waves pass relatively slowly through the asthenosphere compared to the overlying lithospheric mantle. Thus, it has been called the low-velocity zone (LVZ), although the two are not strictly the same; the lower boundary of the LVZ lies at a depth of , whereas the base of the asthenosphere lies at a depth of about . The LVZ also has a high seismic attenuation (seismic waves moving through the asthenosphere lose energy) and significant anisotropy (shear waves polarized vertically have a lower velocity than shear waves polarized horizontally). The discovery of the LVZ alerted seismologists to the existence of the asthenosphere and gave some information about its physical properties, as the speed of seismic waves decreases with decreasing rigidity.
This decrease in seismic wave velocity from the lithosphere to the asthenosphere could be caused by the presence of a very small percentage of melt in the asthenosphere, though since the asthenosphere transmits S waves, it cannot be fully melted.
In the oceanic mantle, the transition from the lithosphere to the asthenosphere (the LAB) is shallower than for the continental mantle (about 60 km in some old oceanic regions) with a sharp and large velocity drop (5–10%). At the mid-ocean ridges, the LAB rises to within a few kilometers of the ocean floor.
The upper part of the asthenosphere is believed to be the zone upon which the great rigid and brittle lithospheric plates of the Earth's crust move about. Due to the temperature and pressure conditions in the asthenosphere, rock becomes ductile, moving at rates of deformation measured in cm/yr over lineal distances eventually measuring thousands of kilometers. In this way, it flows like a convection current, radiating heat outward from the Earth's interior. Above the asthenosphere, at the same rate of deformation, rock behaves elastically and, being brittle, can break, causing faults. The rigid lithosphere is thought to "float" or move about on the slowly flowing asthenosphere, enabling isostatic equilibrium and allowing the movement of tectonic plates.
Boundaries
The asthenosphere extends from an upper boundary at approximately below the surface to a lower boundary at a depth of approximately .
Lithosphere-asthenosphere boundary
The lithosphere-asthenosphere boundary (LAB) is relatively sharp and likely coincides with the onset of partial melting or a change in composition or anisotropy. Various definitions of the boundary reflect various aspects of the boundary region. In addition to the mechanical boundary defined by seismic data, which reflects the transition from the rigid lithosphere to ductile asthenosphere, these include a thermal boundary layer, above which heat is transported by thermal conduction and below which heat transfer is mainly convective; a rheological boundary, where the viscosity drops below about 1021 Pa-s; and a chemical boundary layer, above which the mantle rock is depleted in volatiles and enriched in magnesium relative to the rock below.
Lower boundary of asthenosphere
The lower boundary of the asthenosphere, the top of the tentatively defined mesosphere or mesospheric shell, is less well-defined, but has been placed at the base of the upper mantle. This boundary is neither seismically sharp nor well understood but is approximately coincident with the complex 670 km discontinuity. This discontinuity is generally linked to the transition from mantle rock containing ringwoodite to mantle rock containing bridgmanite and periclase.
Origin
The mechanical properties of the asthenosphere are widely attributed to the partial melting of the rock. It is likely that a small amount of melt is present through much of the asthenosphere, where it is stabilized by the traces of volatiles (water and carbon dioxide) present in the mantle rock. However, the likely amount of melt, not more than about 0.1% of the rock, seems inadequate to fully explain the existence of the asthenosphere. This is not enough melt to fully wet grain boundaries in the rock, and the effects of melt on the mechanical properties of the rock are not expected to be significant if the grain boundaries are not fully wetted. The sharp lithosphere-asthenosphere boundary is also difficult to explain by partial melting alone. It is possible that the asthenosphere is a zone of minimum water solubility in mantle minerals so that more water is available to form greater quantities of melt. Another possible mechanism for producing mechanical weakness is grain boundary sliding, where grains slide slightly past each other under stress, lubricated by the traces of volatiles present. Weakening below oceanic plates is partly caused by their motion itself, thanks to the non-linear dislocation creep mechanism.
Numerical models of mantle convection in which the viscosity is dependent both on temperature and strain rate reliably produce an oceanic asthenosphere, suggesting that strain-rate weakening is a significant contributing mechanism, and explaining the particularly weak asthenosphere below the Pacific plate.
Magma generation
Decompression melting of asthenospheric rock creeping towards the surface is the most important source of magma on Earth. Most of this erupts at mid-ocean ridges to form the distinctive mid-ocean ridge basalt (MORB) of the ocean crust. Magmas are also generated by decompressional melting of the asthenosphere above subduction zones and in areas of continental rifting.
Decompression melting in upwelling asthenosphere likely begins at a depth as great as , where the small amounts of volatiles in the mantle rock (about 100 ppm of water and 60 ppm of carbon dioxide) assist in melting not more than about 0.1% of the rock. At a depth of about , dry melting conditions are reached and melting increases substantially. This dehydrates the remaining solid rock and is likely the origin of the chemically depleted lithosphere.
| Physical sciences | Tectonics | Earth science |
145753 | https://en.wikipedia.org/wiki/Hydrosphere | Hydrosphere | The hydrosphere () is the combined mass of water found on, under, and above the surface of a planet, minor planet, or natural satellite. Although Earth's hydrosphere has been around for about 4 billion years, it continues to change in shape. This is caused by seafloor spreading and continental drift, which rearranges the land and ocean.
It has been estimated that there are 1.386 billion cubic kilometres (333 million cubic miles) of water on Earth. This includes water in gaseous, liquid and frozen forms as soil moisture, groundwater and permafrost in the Earth's crust (to a depth of 2 km); oceans and seas, lakes, rivers and streams, wetlands, glaciers, ice and snow cover on Earth's surface; vapour, droplets and crystals in the air; and part of living plants, animals and unicellular organisms of the biosphere. Saltwater accounts for 97.5% of this amount, whereas fresh water accounts for only 2.5%. Of this fresh water, 68.9% is in the form of ice and permanent snow cover in the Arctic, the Antarctic and mountain glaciers; 30.8% is in the form of fresh groundwater; and only 0.3% of the fresh water on Earth is in easily accessible lakes, reservoirs and river systems.
The total mass of Earth's hydrosphere is about 1.4 × 1018 tonnes, which is about 0.023% of Earth's total mass. At any given time, about 2 × 1013 tonnes of this is in the form of water vapor in the Earth's atmosphere (for practical purposes, 1 cubic metre of water weighs 1 tonne). Approximately 71% of Earth's surface, an area of some 361 million square kilometres (139.5 million square miles), is covered by ocean. The average salinity of Earth's oceans is about 35 grams of salt per kilogram of sea water (3.5%).
History
According to Merriam Webster, the word hydrosphere was brought into English in 1887, translating the German term hydrosphäre, introduced by Eduard Suess.
Water cycle
The water cycle refers to the transfer of water from one state or reservoir to another. Reservoirs include atmospheric moisture (snow, rain and clouds), streams, oceans, rivers, lakes, groundwater, subterranean aquifers, polar ice caps and saturated soil. Solar energy, in the form of heat and light (insolation), and gravity cause the transfer from one state to another over periods from hours to thousands of years. Most evaporation comes from the oceans and is returned to the earth as snow or rain.Sublimation refers to evaporation from snow and ice. Transpiration refers to the expiration of water through the minute pores or stomata of trees. Evapotranspiration is the term used by hydrologists in reference to the three processes together, transpiration, sublimation and evaporation.
Marq de Villiers has described the hydrosphere as a closed system in which water exists. The hydrosphere is intricate, complex, interdependent, all-pervading, stable, and "seems purpose-built for regulating life." De Villiers claimed that, "On earth, the total amount of water has almost certainly not changed since geological times: what we had then we still have. Water can be polluted, abused, and misused but it is neither created nor destroyed, it only migrates. There is no evidence that water vapor escapes into space."Every year the turnover of water on Earth involves 577,000 km3 of water. This is water that evaporates from the oceanic surface (502,800 km3) and from land (74,200 km3). The same amount of water falls as atmospheric precipitation, 458,000 km3 on the ocean and 119,000 km3 on land. The difference between precipitation and evaporation from the land surface (119,000 − 74,200 = 44,800 km3/year) represents the total runoff of the Earth's rivers (42,700 km3/year) and direct groundwater runoff to the ocean (2100 km3/year). These are the principal sources of fresh water to support life necessities and man's economic activities.Water is a basic necessity of life. Since two thirds of the Earth is covered by water, the Earth is also called the blue planet and the watery planet. The hydrosphere plays an important role in the existence of the atmosphere in its present form. Oceans are important in this regard. When the Earth was formed it had only a very thin atmosphere rich in hydrogen and helium similar to the present atmosphere of Mercury. Later the gases hydrogen and helium were expelled from the atmosphere. The gases and water vapor released as the Earth cooled became its present atmosphere. Other gases and water vapor released by volcanoes also entered the atmosphere. As the Earth cooled the water vapor in the atmosphere condensed and fell as rain. The atmosphere cooled further as atmospheric carbon dioxide dissolved into the rain water. In turn, this further caused water vapor to condense and fall as rain. This rain water filled the depressions on the Earth's surface and formed the oceans. It is estimated that this occurred about 4000 million years ago. The first life forms began in the oceans. These organisms did not breathe oxygen. Later, when cyanobacteria evolved, the process of conversion of carbon dioxide into food and oxygen began. As a result, Earth's atmosphere has a distinctly different composition from that of other planets and allowed for life to evolve on Earth.
Human activity has had an impact on the water cycle. Infrastructure, like dams, have a clear, direct impact on the water cycle by blocking and redirecting water pathways. Human caused pollution has changed the biogeochemical cycles of some water systems, and climate change has significantly altered weather patterns. Water withdrawals have exponentially increased because of agriculture, state and domestic use, and infrastructure.
Recharging reservoirs
According to Igor A. Shiklomanov, it takes 2500 years for the complete recharge and replenishment of oceanic waters, 10,000 years for permafrost and ice, 1500 years for deep groundwater and mountainous glaciers, 17 years in lakes, and 16 days in rivers.
Specific fresh water availability
"Specific water availability is the residual (after use) per capita quantity of fresh water." Fresh water resources are unevenly distributed in terms of space and time and can go from floods to water shortages within months in the same area. In 1998, 76% of the total population had a specific water availability of less than 5.0 thousand m3 per year per capita. Already by 1998, 35% of the global population suffered "very low or catastrophically low water supplies," and Shiklomanov predicted that the situation would deteriorate in the twenty-first century with "most of the Earth's population living under the conditions of low or catastrophically low water supply" by 2025. Only 2.5% of the water in the hydrosphere is fresh water and only 0.25% of that water is accessible for our use.
Human impact
The activities of modern humans have drastic effects on the hydrosphere. For instance, water diversion, human development, and pollution all affect the hydrosphere and natural processes within. Humans are withdrawing water from aquifers and diverting rivers at an unprecedented rate. The Ogallala Aquifer is used for agriculture in the United States; if the aquifer goes dry, more than $20 billion worth of food and fiber will vanish from the world's markets. The aquifer is being depleted so much faster than it is replenished that, eventually, the aquifer will run dry. Additionally, only one third of rivers are free-flowing due to the extensive use of dams, levees, hydropower, and habitat degradation. Excessive water use has also caused intermittent streams to become more dry, which is dangerous because they are extremely important for water purification and habitat. Other ways humans impact the hydrosphere include eutrophication, acid rain, and ocean acidification. Humans also rely on the health of the hydrosphere. It is used for water supply, navigation, fishing, agriculture, energy, and recreation.
| Physical sciences | Water: General | null |
145772 | https://en.wikipedia.org/wiki/Food%20web | Food web | A food web is the natural interconnection of food chains and a graphical representation of what-eats-what in an ecological community. Position in the food web, or trophic level, is used in ecology to broadly classify organisms as autotrophs or heterotrophs. This is a non-binary classification; some organisms (such as carnivorous plants) occupy the role of mixotrophs, or autotrophs that additionally obtain organic matter from non-atmospheric sources.
The linkages in a food web illustrate the feeding pathways, such as where heterotrophs obtain organic matter by feeding on autotrophs and other heterotrophs. The food web is a simplified illustration of the various methods of feeding that link an ecosystem into a unified system of exchange. There are different kinds of consumer–resource interactions that can be roughly divided into herbivory, carnivory, scavenging, and parasitism. Some of the organic matter eaten by heterotrophs, such as sugars, provides energy. Autotrophs and heterotrophs come in all sizes, from microscopic to many tonnes - from cyanobacteria to giant redwoods, and from viruses and bdellovibrio to blue whales.
Charles Elton pioneered the concept of food cycles, food chains, and food size in his classical 1927 book "Animal Ecology"; Elton's 'food cycle' was replaced by 'food web' in a subsequent ecological text. Elton organized species into functional groups, which was the basis for Raymond Lindeman's classic and landmark paper in 1942 on trophic dynamics. Lindeman emphasized the important role of decomposer organisms in a trophic system of classification. The notion of a food web has a historical foothold in the writings of Charles Darwin and his terminology, including an "entangled bank", "web of life", "web of complex relations", and in reference to the decomposition actions of earthworms he talked about "the continued movement of the particles of earth". Even earlier, in 1768 John Bruckner described nature as "one continued web of life".
Food webs are limited representations of real ecosystems as they necessarily aggregate many species into trophic species, which are functional groups of species that have the same predators and prey in a food web. Ecologists use these simplifications in quantitative (or mathematical representation) models of trophic or consumer-resource systems dynamics. Using these models they can measure and test for generalized patterns in the structure of real food web networks. Ecologists have identified non-random properties in the topological structure of food webs. Published examples that are used in meta analysis are of variable quality with omissions. However, the number of empirical studies on community webs is on the rise and the mathematical treatment of food webs using network theory had identified patterns that are common to all. Scaling laws, for example, predict a relationship between the topology of food web predator-prey linkages and levels of species richness.
Taxonomy of a food web
Links in food webs map the feeding connections (who eats whom) in an ecological community. Food cycle is an obsolete term that is synonymous with food web. Ecologists can broadly group all life forms into one of two trophic layers, the autotrophs and the heterotrophs. Autotrophs produce more biomass energy, either chemically without the sun's energy or by capturing the sun's energy in photosynthesis, than they use during metabolic respiration. Heterotrophs consume rather than produce biomass energy as they metabolize, grow, and add to levels of secondary production. A food web depicts a collection of polyphagous heterotrophic consumers that network and cycle the flow of energy and nutrients from a productive base of self-feeding autotrophs.
The base or basal species in a food web are those species without prey and can include autotrophs or saprophytic detritivores (i.e., the community of decomposers in soil, biofilms, and periphyton). Feeding connections in the web are called trophic links. The number of trophic links per consumer is a measure of food web connectance. Food chains are nested within the trophic links of food webs. Food chains are linear (noncyclic) feeding pathways that trace monophagous consumers from a base species up to the top consumer, which is usually a larger predatory carnivore.
Linkages connect to nodes in a food web, which are aggregates of biological taxa called trophic species. Trophic species are functional groups that have the same predators and prey in a food web. Common examples of an aggregated node in a food web might include parasites, microbes, decomposers, saprotrophs, consumers, or predators, each containing many species in a web that can otherwise be connected to other trophic species.
Trophic levels
Food webs have trophic levels and positions. Basal species, such as plants, form the first level and are the resource-limited species that feed on no other living creature in the web. Basal species can be autotrophs or detritivores, including "decomposing organic material and its associated microorganisms which we defined as detritus, micro-inorganic material and associated microorganisms (MIP), and vascular plant material." Most autotrophs capture the sun's energy in chlorophyll, but some autotrophs (the chemolithotrophs) obtain energy by the chemical oxidation of inorganic compounds and can grow in dark environments, such as the sulfur bacterium Thiobacillus, which lives in hot sulfur springs. The top level has top (or apex) predators that no other species kills directly for their food resource needs. The intermediate levels are filled with omnivores that feed on more than one trophic level and cause energy to flow through several food pathways starting from a basal species.
In the simplest scheme, the first trophic level (level 1) is plants, then herbivores (level 2), and then carnivores (level 3). The trophic level equals one more than the chain length, which is the number of links connecting to the base. The base of the food chain (primary producers or detritivores) is set at zero. Ecologists identify feeding relations and organize species into trophic species through extensive gut content analysis of different species. The technique has been improved through the use of stable isotopes to better trace energy flow through the web. It was once thought that omnivory was rare, but recent evidence suggests otherwise. This realization has made trophic classifications more complex.
Trophic dynamics and multitrophic interactions
The trophic level concept was introduced in a historical landmark paper on trophic dynamics in 1942 by Raymond L. Lindeman. The basis of trophic dynamics is the transfer of energy from one part of the ecosystem to another. The trophic dynamic concept has served as a useful quantitative heuristic, but it has several major limitations including the precision by which an organism can be allocated to a specific trophic level. Omnivores, for example, are not restricted to any single level. Nonetheless, recent research has found that discrete trophic levels do exist, but "above the herbivore trophic level, food webs are better characterized as a tangled web of omnivores."
A central question in the trophic dynamic literature is the nature of control and regulation over resources and production. Ecologists use simplified one trophic position food chain models (producer, carnivore, decomposer). Using these models, ecologists have tested various types of ecological control mechanisms. For example, herbivores generally have an abundance of vegetative resources, which meant that their populations were largely controlled or regulated by predators. This is known as the top-down hypothesis or 'green-world' hypothesis. Alternatively to the top-down hypothesis, not all plant material is edible and the nutritional quality or antiherbivore defenses of plants (structural and chemical) suggests a bottom-up form of regulation or control. Recent studies have concluded that both "top-down" and "bottom-up" forces can influence community structure and the strength of the influence is environmentally context dependent. These complex multitrophic interactions involve more than two trophic levels in a food web. For example, such interactions have been discovered in the context of arbuscular mycorrhizal fungi and aphid herbivores that utilize the same plant species.
Another example of a multitrophic interaction is a trophic cascade, in which predators help to increase plant growth and prevent overgrazing by suppressing herbivores. Links in a food-web illustrate direct trophic relations among species, but there are also indirect effects that can alter the abundance, distribution, or biomass in the trophic levels. For example, predators eating herbivores indirectly influence the control and regulation of primary production in plants. Although the predators do not eat the plants directly, they regulate the population of herbivores that are directly linked to plant trophism. The net effect of direct and indirect relations is called trophic cascades. Trophic cascades are separated into species-level cascades, where only a subset of the food-web dynamic is impacted by a change in population numbers, and community-level cascades, where a change in population numbers has a dramatic effect on the entire food-web, such as the distribution of plant biomass.
The field of chemical ecology has elucidated multitrophic interactions that entail the transfer of defensive compounds across multiple trophic levels. For example, certain plant species in the Castilleja and Plantago genera have been found to produce defensive compounds called iridoid glycosides that are sequestered in the tissues of the Taylor's checkerspot butterfly larvae that have developed a tolerance for these compounds and are able to consume the foliage of these plants. These sequestered iridoid glycosides then confer chemical protection against bird predators to the butterfly larvae. Another example of this sort of multitrophic interaction in plants is the transfer of defensive alkaloids produced by endophytes living within a grass host to a hemiparasitic plant that is also using the grass as a host.
Energy flow and biomass
Food webs depict energy flow via trophic linkages. Energy flow is directional, which contrasts against the cyclic flows of material through the food web systems. Energy flow "typically includes production, consumption, assimilation, non-assimilation losses (feces), and respiration (maintenance costs)." In a very general sense, energy flow (E) can be defined as the sum of metabolic production (P) and respiration (R), such that E=P+R.
Biomass represents stored energy. However, concentration and quality of nutrients and energy is variable. Many plant fibers, for example, are indigestible to many herbivores leaving grazer community food webs more nutrient limited than detrital food webs where bacteria are able to access and release the nutrient and energy stores. "Organisms usually extract energy in the form of carbohydrates, lipids, and proteins. These polymers have a dual role as supplies of energy as well as building blocks; the part that functions as energy supply results in the production of nutrients (and carbon dioxide, water, and heat). Excretion of nutrients is, therefore, basic to metabolism." The units in energy flow webs are typically a measure mass or energy per m2 per unit time. Different consumers are going to have different metabolic assimilation efficiencies in their diets. Each trophic level transforms energy into biomass. Energy flow diagrams illustrate the rates and efficiency of transfer from one trophic level into another and up through the hierarchy.
It is the case that the biomass of each trophic level decreases from the base of the chain to the top. This is because energy is lost to the environment with each transfer as entropy increases. About eighty to ninety percent of the energy is expended for the organism's life processes or is lost as heat or waste. Only about ten to twenty percent of the organism's energy is generally passed to the next organism. The amount can be less than one percent in animals consuming less digestible plants, and it can be as high as forty percent in zooplankton consuming phytoplankton. Graphic representations of the biomass or productivity at each tropic level are called ecological pyramids or trophic pyramids. The transfer of energy from primary producers to top consumers can also be characterized by energy flow diagrams.
Food chain
A common metric used to quantify food web trophic structure is food chain length. Food chain length is another way of describing food webs as a measure of the number of species encountered as energy or nutrients move from the plants to top predators. There are different ways of calculating food chain length depending on what parameters of the food web dynamic are being considered: connectance, energy, or interaction. In its simplest form, the length of a chain is the number of links between a trophic consumer and the base of the web. The mean chain length of an entire web is the arithmetic average of the lengths of all chains in a food web.
In a simple predator-prey example, a deer is one step removed from the plants it eats (chain length = 1) and a wolf that eats the deer is two steps removed from the plants (chain length = 2). The relative amount or strength of influence that these parameters have on the food web address questions about:
the identity or existence of a few dominant species (called strong interactors or keystone species)
the total number of species and food-chain length (including many weak interactors) and
how community structure, function and stability is determined.
Ecological pyramids
In a pyramid of numbers, the number of consumers at each level decreases significantly, so that a single top consumer, (e.g., a polar bear or a human), will be supported by a much larger number of separate producers. There is usually a maximum of four or five links in a food chain, although food chains in aquatic ecosystems are more often longer than those on land. Eventually, all the energy in a food chain is dispersed as heat.
Ecological pyramids place the primary producers at the base. They can depict different numerical properties of ecosystems, including numbers of individuals per unit of area, biomass (g/m2), and energy (k cal m−2 yr−1). The emergent pyramidal arrangement of trophic levels with amounts of energy transfer decreasing as species become further removed from the source of production is one of several patterns that is repeated amongst the planets ecosystems. The size of each level in the pyramid generally represents biomass, which can be measured as the dry weight of an organism. Autotrophs may have the highest global proportion of biomass, but they are closely rivaled or surpassed by microbes.
Pyramid structure can vary across ecosystems and across time. In some instances biomass pyramids can be inverted. This pattern is often identified in aquatic and coral reef ecosystems. The pattern of biomass inversion is attributed to different sizes of producers. Aquatic communities are often dominated by producers that are smaller than the consumers that have high growth rates. Aquatic producers, such as planktonic algae or aquatic plants, lack the large accumulation of secondary growth as exists in the woody trees of terrestrial ecosystems. However, they are able to reproduce quickly enough to support a larger biomass of grazers. This inverts the pyramid. Primary consumers have longer lifespans and slower growth rates that accumulates more biomass than the producers they consume. Phytoplankton live just a few days, whereas the zooplankton eating the phytoplankton live for several weeks and the fish eating the zooplankton live for several consecutive years. Aquatic predators also tend to have a lower death rate than the smaller consumers, which contributes to the inverted pyramidal pattern. Population structure, migration rates, and environmental refuge for prey are other possible causes for pyramids with biomass inverted. Energy pyramids, however, will always have an upright pyramid shape if all sources of food energy are included and this is dictated by the second law of thermodynamics.
Material flux and recycling
Many of the Earth's elements and minerals (or mineral nutrients) are contained within the tissues and diets of organisms. Hence, mineral and nutrient cycles trace food web energy pathways. Ecologists employ stoichiometry to analyze the ratios of the main elements found in all organisms: carbon (C), nitrogen (N), phosphorus (P). There is a large transitional difference between many terrestrial and aquatic systems as C:P and C:N ratios are much higher in terrestrial systems while N:P ratios are equal between the two systems. Mineral nutrients are the material resources that organisms need for growth, development, and vitality. Food webs depict the pathways of mineral nutrient cycling as they flow through organisms. Most of the primary production in an ecosystem is not consumed, but is recycled by detritus back into useful nutrients. Many of the Earth's microorganisms are involved in the formation of minerals in a process called biomineralization. Bacteria that live in detrital sediments create and cycle nutrients and biominerals. Food web models and nutrient cycles have traditionally been treated separately, but there is a strong functional connection between the two in terms of stability, flux, sources, sinks, and recycling of mineral nutrients.
Kinds of food webs
Food webs are necessarily aggregated and only illustrate a tiny portion of the complexity of real ecosystems. For example, the number of species on the planet are likely in the general order of 107, over 95% of these species consist of microbes and invertebrates, and relatively few have been named or classified by taxonomists. It is explicitly understood that natural systems are 'sloppy' and that food web trophic positions simplify the complexity of real systems that sometimes overemphasize many rare interactions. Most studies focus on the larger influences where the bulk of energy transfer occurs. "These omissions and problems are causes for concern, but on present evidence do not present insurmountable difficulties."
There are different kinds or categories of food webs:
Source web - one or more node(s), all of their predators, all the food these predators eat, and so on.
Sink web - one or more node(s), all of their prey, all the food that these prey eat, and so on.
Community (or connectedness) web - a group of nodes and all the connections of who eats whom.
Energy flow web - quantified fluxes of energy between nodes along links between a resource and a consumer.
Paleoecological web - a web that reconstructs ecosystems from the fossil record.
Functional web - emphasizes the functional significance of certain connections having strong interaction strength and greater bearing on community organization, more so than energy flow pathways. Functional webs have compartments, which are sub-groups in the larger network where there are different densities and strengths of interaction. Functional webs emphasize that "the importance of each population in maintaining the integrity of a community is reflected in its influence on the growth rates of other populations."
Within these categories, food webs can be further organized according to the different kinds of ecosystems being investigated. For example, human food webs, agricultural food webs, detrital food webs, marine food webs, aquatic food webs, soil food webs, Arctic (or polar) food webs, terrestrial food webs, and microbial food webs. These characterizations stem from the ecosystem concept, which assumes that the phenomena under investigation (interactions and feedback loops) are sufficient to explain patterns within boundaries, such as the edge of a forest, an island, a shoreline, or some other pronounced physical characteristic.
Detrital web
In a detrital web, plant and animal matter is broken down by decomposers, e.g., bacteria and fungi, and moves to detritivores and then carnivores. There are often relationships between the detrital web and the grazing web. Mushrooms produced by decomposers in the detrital web become a food source for deer, squirrels, and mice in the grazing web. Earthworms eaten by robins are detritivores consuming decaying leaves.
"Detritus can be broadly defined as any form of non-living organic matter, including different types of plant tissue (e.g. leaf litter, dead wood, aquatic macrophytes, algae), animal tissue (carrion), dead microbes, faeces (manure, dung, faecal pellets, guano, frass), as well as products secreted, excreted or exuded from organisms (e.g. extra-cellular polymers, nectar, root exudates and leachates, dissolved organic matter, extra-cellular matrix, mucilage). The relative importance of these forms of detritus, in terms of origin, size and chemical composition, varies across ecosystems."
Quantitative food webs
Ecologists collect data on trophic levels and food webs to statistically model and mathematically calculate parameters, such as those used in other kinds of network analysis (e.g., graph theory), to study emergent patterns and properties shared among ecosystems. There are different ecological dimensions that can be mapped to create more complicated food webs, including: species composition (type of species), richness (number of species), biomass (the dry weight of plants and animals), productivity (rates of conversion of energy and nutrients into growth), and stability (food webs over time). A food web diagram illustrating species composition shows how change in a single species can directly and indirectly influence many others. Microcosm studies are used to simplify food web research into semi-isolated units such as small springs, decaying logs, and laboratory experiments using organisms that reproduce quickly, such as daphnia feeding on algae grown under controlled environments in jars of water.
While the complexity of real food webs connections are difficult to decipher, ecologists have found mathematical models on networks an invaluable tool for gaining insight into the structure, stability, and laws of food web behaviours relative to observable outcomes. "Food web theory centers around the idea of connectance." Quantitative formulas simplify the complexity of food web structure. The number of trophic links (tL), for example, is converted into a connectance value:
,
where, S(S-1)/2 is the maximum number of binary connections among S species. "Connectance (C) is the fraction of all possible links that are realized (L/S2) and represents a standard measure of food web complexity..." The distance (d) between every species pair in a web is averaged to compute the mean distance between all nodes in a web (D) and multiplied by the total number of links (L) to obtain link-density (LD), which is influenced by scale-dependent variables such as species richness. These formulas are the basis for comparing and investigating the nature of non-random patterns in the structure of food web networks among many different types of ecosystems.
Scaling laws, complexity, chaos, and pattern correlates are common features attributed to food web structure.
Complexity and stability
Food webs are extremely complex. Complexity is a term that conveys the mental intractability of understanding all possible higher-order effects in a food web. Sometimes in food web terminology, complexity is defined as product of the number of species and connectance., though there have been criticisms of this definition and other proposed methods for measuring network complexity. Connectance is "the fraction of all possible links that are realized in a network". These concepts were derived and stimulated through the suggestion that complexity leads to stability in food webs, such as increasing the number of trophic levels in more species rich ecosystems. This hypothesis was challenged through mathematical models suggesting otherwise, but subsequent studies have shown that the premise holds in real systems.
At different levels in the hierarchy of life, such as the stability of a food web, "the same overall structure is maintained in spite of an ongoing flow and change of components." The farther a living system (e.g., ecosystem) sways from equilibrium, the greater its complexity. Complexity has multiple meanings in the life sciences and in the public sphere that confuse its application as a precise term for analytical purposes in science. Complexity in the life sciences (or biocomplexity) is defined by the "properties emerging from the interplay of behavioral, biological, physical, and social interactions that affect, sustain, or are modified by living organisms, including humans".
Several concepts have emerged from the study of complexity in food webs. Complexity explains many principals pertaining to self-organization, non-linearity, interaction, cybernetic feedback, discontinuity, emergence, and stability in food webs. Nestedness, for example, is defined as "a pattern of interaction in which specialists interact with species that form perfect subsets of the species with which generalists interact", "—that is, the diet of the most specialized species is a subset of the diet of the next more generalized species, and its diet a subset of the next more generalized, and so on." Until recently, it was thought that food webs had little nested structure, but empirical evidence shows that many published webs have nested subwebs in their assembly.
Food webs are complex networks. As networks, they exhibit similar structural properties and mathematical laws that have been used to describe other complex systems, such as small world and scale free properties. The small world attribute refers to the many loosely connected nodes, non-random dense clustering of a few nodes (i.e., trophic or keystone species in ecology), and small path length compared to a regular lattice. "Ecological networks, especially mutualistic networks, are generally very heterogeneous, consisting of areas with sparse links among species and distinct areas of tightly linked species. These regions of high link density are often referred to as cliques, hubs, compartments, cohesive sub-groups, or modules...Within food webs, especially in aquatic systems, nestedness appears to be related to body size because the diets of smaller predators tend to be nested subsets of those of larger predators (Woodward & Warren 2007; YvonDurocher et al. 2008), and phylogenetic constraints, whereby related taxa are nested based on their common evolutionary history, are also evident (Cattin et al. 2004)." "Compartments in food webs are subgroups of taxa in which many strong interactions occur within the subgroups and few weak interactions occur between the subgroups. Theoretically, compartments increase the stability in networks, such as food webs."
Food webs are also complex in the way that they change in scale, seasonally, and geographically. The components of food webs, including organisms and mineral nutrients, cross the thresholds of ecosystem boundaries. This has led to the concept or area of study known as cross-boundary subsidy. "This leads to anomalies, such as food web calculations determining that an ecosystem can support one half of a top carnivore, without specifying which end." Nonetheless, real differences in structure and function have been identified when comparing different kinds of ecological food webs, such as terrestrial vs. aquatic food webs.
History of food webs
Food webs serve as a framework to help ecologists organize the complex network of interactions among species observed in nature and around the world. One of the earliest descriptions of a food chain was described by a medieval Afro-Arab scholar named Al-Jahiz: "All animals, in short, cannot exist without food, neither can the hunting animal escape being hunted in his turn." The earliest graphical depiction of a food web was by Lorenzo Camerano in 1880, followed independently by those of Pierce and colleagues in 1912 and Victor Shelford in 1913. Two food webs about herring were produced by Victor Summerhayes and Charles Elton and Alister Hardy in 1923 and 1924. Charles Elton subsequently pioneered the concept of food cycles, food chains, and food size in his classical 1927 book "Animal Ecology"; Elton's 'food cycle' was replaced by 'food web' in a subsequent ecological text. After Charles Elton's use of food webs in his 1927 synthesis, they became a central concept in the field of ecology. Elton organized species into functional groups, which formed the basis for the trophic system of classification in Raymond Lindeman's classic and landmark paper in 1942 on trophic dynamics. The notion of a food web has a historical foothold in the writings of Charles Darwin and his terminology, including an "entangled bank", "web of life", "web of complex relations", and in reference to the decomposition actions of earthworms he talked about "the continued movement of the particles of earth". Even earlier, in 1768 John Bruckner described nature as "one continued web of life".
Interest in food webs increased after Robert Paine's experimental and descriptive study of intertidal shores suggesting that food web complexity was key to maintaining species diversity and ecological stability. Many theoretical ecologists, including Sir Robert May and Stuart Pimm, were prompted by this discovery and others to examine the mathematical properties of food webs.
| Biology and health sciences | Ecology | Biology |
145773 | https://en.wikipedia.org/wiki/Buoy | Buoy | A buoy (; ) is a floating device that can have many purposes. It can be anchored (stationary) or allowed to drift with ocean currents.
History
The ultimate origin of buoys is unknown, but by 1295 a seaman's manual referred to navigation buoys in the Guadalquivir River in Spain. To the north there are early medieval mentions of the French / Belgian River Maas being buoyed. Such early buoys were probably just timber beams or rafts, but in 1358 there is a record of a barrel buoy in the Dutch Maasmond (also known as the Maas Sluis or Maasgat). The simple barrel was difficult to secure to the seabed, and so a conical tonne was developed. They had a solid plug at the narrow end through which a mooring ring could be attached. By 1790 the older conical tonne was being replaced by a nun buoy. This had the same conical section below the waterline as the tonne buoy, but at the waterline a barrel shape was used to allow a truncated cone to be above the water. The whole was completed with a top mark. In the nineteenth century iron buoys became available. They had watertight internal bulkheads and as well as topmarks and might have bells (1860) or whistles (1880). In 1879 Julius Pintsch obtained a patent for the illumination of buoys by using a compressed gas. This was superseded from 1912 onwards by Gustaf Dalén's acetylene lamp. This could be set to flash which ensured that buoys could be distinguished from ships' lights and from each other. A later development was the sun valve which shut off the gas during sunlight.
Types
Navigational buoys
Race course marker buoys are used for buoy racing, the most prevalent form of yacht racing and power boat racing. They delimit the course and must be passed to a specified side. They are also used in underwater orienteering competitions.
Emergency wreck buoys provide a clear and unambiguous means of temporarily marking new wrecks, typically for the first 24–72 hours. They are coloured in an equal number of blue and yellow vertical stripes and fitted with an alternating blue and yellow flashing light. They were implemented following collisions in the Dover Strait in 2002 when vessels struck the new wreck of the .
Ice marking buoys mark holes in frozen lakes and rivers so snowmobiles do not drive over the holes.
Large Navigational Buoys (LNB, or Lanby buoys) are automatic buoys over 10 meters high equipped with a powerful light monitored electronically as a replacement for a lightvessel. They may be marked on charts as a "Superbuoy."
Lateral marker buoys
Safe water mark or fairway buoys mark the entrance to a channel or nearby landfall
Sea marks aid pilotage by marking a maritime channel, hazard or administrative area to allow boats and ships to navigate safely. Some are fitted with wave-activated bells or gongs.
Wreck buoys mark a wrecked ship to warn other ships to keep away because of unseen hazards.
Light buoys provide demarcation at night.
Marker buoys
Buoys are often used to temporarily or permanently mark the positions of underwater objects:
Lobster trap buoys are brightly colored buoys marking lobster trap locations so lobster fishers can find their lobster traps. Each fisher has a unique colour marking or registration number. They are allowed to haul only their own traps, and must display their buoy colour or license number on their boat so law enforcement officials know what they should be hauling. The buoys are brightly coloured with highly visible numbers so they can be seen in poor visibility conditions like rain, fog and sea smoke.
Fishing floats are a type of lightweight buoys used in angling to mark the position of the baited hook suspended underneath, and as a bite indicator to signal the angler any changes in the hook's underwater status.
Diving
Several types of marker buoys may be used by divers:
Decompression buoys are deployed by submerged SCUBA divers to mark their position underwater whilst doing decompression stops
Shot buoys mark dive sites for the boat safety cover of scuba divers so they can descend to dive sites more easily in conditions of low visibility or tidal currents and more safely do decompression stops on their ascents.
Surface marker buoys are taken on dives by scuba divers to mark their positions underwater.
Dive site demarcation buoys indicate that divers are working in the marked area, to warn passing vessels to stay clear.
Rescue
Lifebuoys are lifesaving buoys thrown to people in the water to provide buoyancy. They usually have a connected line allowing them to be pulled in.
Self-locating datum marker buoys (SLDMB) are 70% scale Coastal Ocean Dynamics Experiment (CODE)/Davis-style oceanographic surface drifters with drogue vanes between 30 and 100 cm deep, designed for deployment from U.S. Coast Guard vessels or airframes for search and rescue. They have very little surface area above water to minimize the effects on them off winds and waves.
Submarine rescue buoys are released in emergencies and for communication purposes.
Research
Profiling buoys are specialized buoys that adjust their buoyancy to sink at a controlled rate to a set depth, for example 2,000 metres while measuring sea temperature and salinity. After a certain period, typically 10 days, they return to the surface, transmit their data via satellite, then sink again. See Argo (oceanography).
Tsunami buoys are anchored buoys that can detect sudden changes in undersea water pressure, and are a component of tsunami warning systems in the Pacific Tsunami Warning Center and Indian Oceans.
Wave buoys measure the movement of the water surface as a wave train. The data they transmit is analysed to form statistics like significant wave height and period, and wave direction.
Weather buoys measure weather parameters such as air temperature, barometric pressure, and wind speed and direction. They transmit this data, via satellite radio links such as the purpose-built Argos System or commercial satellite phone networks, to meteorological centres for forecasting and climate study. They may be anchored (moored buoys), or allowed to drift (drifting buoys) in the open currents. Their position is calculated by the satellite. They are also referred to as Ocean Data Acquisition Systems, or (ODAS) buoys. and may be marked on charts as "Superbuoys."
Mooring
Mooring buoys keep one end of a mooring cable or chain on the water's surface so ships and boats can tie to them. Many marinas mark them with numbers and assign them to particular vessels, or rent them to transient vessels.
Tripping buoys are used to keep one end of a '' to be used to break out and lift an anchor on the water's surface so that a stuck anchor can more easily be freed.
Military
Marker buoys, used in naval warfare (particularly anti-submarine warfare) emit light and/or smoke using pyrotechnic devices to create the flare and smoke. Commonly 3 inches (76 mm) in diameter and about 20 inches (500 mm) long, they are activated by contact with seawater and float on the surface. Some extinguish themselves after a specific period, while others are sunk when they are no longer needed.
Sonobuoys are used by anti-submarine warfare aircraft to detect submarines by SONAR.
Target buoys simulate targets, such as small boats, in live-fire exercises by naval and coastal forces. They are usually targeted by medium-sized weapons such as heavy machine guns, rapid fire cannons (~20 mm), autocannons (up to 40–57 mm) and anti-tank rockets.
Specific forms
DAN buoys are used as:
Large maritime navigational aids providing a platform for light and radio beacons
Lifebuoys with flags, used on yachts and smaller pleasure craft
Temporary markers in Danish seine fishing to mark net anchor positions
Temporary markers set by danlayers during minesweeping operations to indicate the boundaries of swept paths, swept areas, known hazards, and other locations or reference points
Temporary markers for rescue operations
Spar buoys are tall, thin buoys that float upright, e.g. R/P FLIP
Other
Letter boxes on buoys exist in Töre (Sweden) and at the Steinhuder Meer (Germany)
Fictional
Imaginary "Mail buoys" have been used as a prank in the US Navy when a new sailor may be given the task of locating one to retrieve non-existent mail.
Space buoys, a feature in some science fiction stories which are stationary objects in outer space that provide navigation data or warnings.
Other uses
The word "buoyed" can also be used figuratively. For example, a person can buoy up ('lift up') someone's spirits by providing help and empathy.
Buoys are used in some wave power systems to generate electrical power.
George A. Stephen, founder of Weber-Stephen Products Co., invented the kettle grill by cutting a metal buoy in half and fashioning a dome shaped grill to it with a rounded lid.
Gallery
| Technology | Naval transport | null |
145813 | https://en.wikipedia.org/wiki/Subduction | Subduction | Subduction is a geological process in which the oceanic lithosphere and some continental lithosphere is recycled into the Earth's mantle at the convergent boundaries between tectonic plates. Where one tectonic plate converges with a second plate, the heavier plate dives beneath the other and sinks into the mantle. A region where this process occurs is known as a subduction zone, and its surface expression is known as an arc-trench complex. The process of subduction has created most of the Earth's continental crust. Rates of subduction are typically measured in centimeters per year, with rates of convergence as high as 11 cm/year.
Subduction is possible because the cold and rigid oceanic lithosphere is slightly denser than the underlying asthenosphere, the hot, ductile layer in the upper mantle. Once initiated, stable subduction is driven mostly by the negative buoyancy of the dense subducting lithosphere. The down-going slab sinks into the mantle largely under its own weight.
Earthquakes are common along subduction zones, and fluids released by the subducting plate trigger volcanism in the overriding plate. If the subducting plate sinks at a shallow angle, the overriding plate develops a belt of deformation characterized by crustal thickening, mountain building, and metamorphism. Subduction at a steeper angle is characterized by the formation of back-arc basins.
Subduction and plate tectonics
According to the theory of plate tectonics, the Earth's lithosphere, its rigid outer shell, is broken into sixteen larger tectonic plates and several smaller plates. These plates are in slow motion, due mostly to the pull force of subducting lithosphere. Sinking lithosphere at subduction zones are a part of convection cells in the underlying ductile mantle. This process of convection allows heat generated by radioactive decay to escape from the Earth's interior.
The lithosphere consists of the outermost light crust plus the uppermost rigid portion of the mantle. Oceanic lithosphere ranges in thickness from just a few km for young lithosphere created at mid-ocean ridges to around for the oldest oceanic lithosphere. Continental lithosphere is up to thick. The lithosphere is relatively cold and rigid compared with the underlying asthenosphere, and so tectonic plates move as solid bodies atop the asthenosphere. Individual plates often include both regions of the oceanic lithosphere and continental lithosphere.
Subduction zones are where cold oceanic lithosphere sinks back into the mantle and is recycled. They are found at convergent plate boundaries, where the heavier oceanic lithosphere of one plate is overridden by the leading edge of another, less-dense plate. The overridden plate (the slab) sinks at an angle most commonly between 25 and 75 degrees to Earth's surface. This sinking is driven by the temperature difference between the slab and the surrounding asthenosphere, as the colder oceanic lithosphere is, on average, more dense. Sediments and some trapped water are carried downwards by the slab and recycled into the deep mantle.
So far, Earth is the only planet where subduction is known to occur, and subduction zones are its most important tectonic feature. Subduction is the driving force behind plate tectonics, and without it, plate tectonics could not occur. Oceanic subduction zones are located along of convergent plate margins, almost equal to the cumulative plate formation rate of mid-ocean ridges.
Sea water seeps into oceanic lithosphere through fractures and pores, and reacts with minerals in the crust and mantle to form hydrous minerals (such as serpentine) that store water in their crystal structures. Water is transported into the deep mantle via hydrous minerals in subducting slabs. During subduction, a series of minerals in these slabs such as serpentine can be stable at different pressures within the slab geotherms, and may transport significant amount of water into the Earth's interior. As plates sink and heat up, released fluids can trigger seismicity and induce melting within the subducted plate and in the overlying mantle wedge. This type of melting selectively concentrates volatiles and transports them into the overlying plate. If an eruption occurs, the cycle then returns the volatiles into the oceans and atmosphere.
Structure of subduction zones
Arc-trench complex
The surface expressions of subduction zones are arc-trench complexes. On the ocean side of the complex, where the subducting plate first approaches the subduction zone, there is often an outer trench high or outer trench swell. Here the plate shallows slightly before plunging downwards, as a consequence of the rigidity of the plate. The point where the slab begins to plunge downwards is marked by an oceanic trench. Oceanic trenches are the deepest parts of the ocean floor.
Beyond the trench is the forearc portion of the overriding plate. Depending on sedimentation rates, the forearc may include an accretionary wedge of sediments scraped off the subducting slab and accreted to the overriding plate. However, not all arc-trench complexes have an accretionary wedge. Accretionary arcs have a well-developed forearc basin behind the accretionary wedge, while the forearc basin is poorly developed in non-accretionary arcs.
Beyond the forearc basin, volcanoes are found in long chains called volcanic arcs. The subducting basalt and sediment are normally rich in hydrous minerals and clays. Additionally, large quantities of water are introduced into cracks and fractures created as the subducting slab bends downward. During the transition from basalt to eclogite, these hydrous materials break down, producing copious quantities of water, which at such great pressure and temperature exists as a supercritical fluid. The supercritical water, which is hot and more buoyant than the surrounding rock, rises into the overlying mantle, where it lowers the melting temperature of the mantle rock, generating magma via flux melting. The magmas, in turn, rise as diapirs because they are less dense than the rocks of the mantle. The mantle-derived magmas (which are initially basaltic in composition) can ultimately reach the Earth's surface, resulting in volcanic eruptions. The chemical composition of the erupting lava depends upon the degree to which the mantle-derived basalt interacts with (melts) Earth's crust or undergoes fractional crystallization. Arc volcanoes tend to produce dangerous eruptions because they are rich in water (from the slab and sediments) and tend to be extremely explosive. Krakatoa, Nevado del Ruiz, and Mount Vesuvius are all examples of arc volcanoes. Arcs are also associated with most ore deposits.
Beyond the volcanic arc is a back-arc region whose character depends strongly on the angle of subduction of the subducting slab. Where this angle is shallow, the subducting slab drags the overlying continental crust partially with it, which produces a zone of shortening and crustal thickening in which there may be extensive folding and thrust faulting. If the angle of subduction steepens or rolls back, the upper plate lithosphere will be put in tension instead, often producing a back-arc basin.
Deep structure
The arc-trench complex is the surface expression of a much deeper structure. Though not directly accessible, the deeper portions can be studied using geophysics and geochemistry. Subduction zones are defined by an inclined zone of earthquakes, the Wadati–Benioff zone, that dips away from the trench and extends down below the volcanic arc to the 660-kilometer discontinuity. Subduction zone earthquakes occur at greater depths (up to ) than elsewhere on Earth (typically less than depth); such deep earthquakes may be driven by deep phase transformations, thermal runaway, or dehydration embrittlement. Seismic tomography shows that some slabs can penetrate the lower mantle and sink clear to the core–mantle boundary. Here the residue of the slabs may eventually heat enough to rise back to the surface as mantle plumes.
Subduction angle
Subduction typically occurs at a moderately steep angle by the time it is beneath the volcanic arc. However, anomalous shallower angles of subduction are known to exist as well as some that are extremely steep.
Flat slab subduction (subducting angle less than 30°) occurs when the slab subducts nearly horizontally. The relatively flat slab can extend for hundreds of kilometers under the upper plate. This geometry is commonly caused by the subduction of buoyant lithosphere due to thickened crust or warmer lithosphere. Recent studies have also shown a strong correlation that older and wider subduction zones are related to flatter subduction dips. This provides an explanation as to why flat subduction only presently occur in the eastern pacific as only these regions were old and wide enough to support flat slab subduction and why the Laramide flat slab subduction and South China flat slab subduction were possible. Hu ultimately proposes that a combination of subduction age and slab characteristics provide the strongest controls over subduction dips. Because subduction of slabs to depth is necessary to drive subduction zone volcanism, flat-slab subduction can be invoked to explain volcanic gaps.
Flat-slab subduction is ongoing beneath part of the Andes, causing segmentation of the Andean Volcanic Belt into four zones. The flat-slab subduction in northern Peru and the Norte Chico region of Chile is believed to be the result of the subduction of two buoyant aseismic ridges, the Nazca Ridge and the Juan Fernández Ridge, respectively. Around Taitao Peninsula flat-slab subduction is attributed to the subduction of the Chile Rise, a spreading ridge.
The Laramide Orogeny in the Rocky Mountains of the United States is attributed to flat-slab subduction. During this orogeny, a broad volcanic gap appeared at the southwestern margin of North America, and deformation occurred much farther inland; it was during this time that the basement-cored mountain ranges of Colorado, Utah, Wyoming, South Dakota, and New Mexico came into being. The most massive subduction zone earthquakes, so-called "megaquakes", have been found to occur in flat-slab subduction zones.
Steep-angle subduction (subducting angle greater than 70°) occurs in subduction zones where Earth's oceanic crust and lithosphere are cold and thick and have, therefore, lost buoyancy. Recent studies have also correlated steep angled subduction zones with younger and less extensive subduction zones. This would explain why most modern subduction zones are relatively steep. The steepest dipping subduction zone lies in the Mariana Trench, which is also where the oceanic lithosphere of Jurassic age is the oldest on Earth exempting ophiolites. Steep-angle subduction is, in contrast to flat-slab subduction, associated with back-arc extension of the upper plate, creating volcanic arcs and pulling fragments of continental crust away from continents to leave behind a marginal sea.
Life cycle of subduction zones
Initiation of subduction
Although stable subduction is fairly well understood, the process by which subduction is initiated remains a matter of discussion and continuing study. Subduction can begin spontaneously if the denser oceanic lithosphere can founder and sink beneath the adjacent oceanic or continental lithosphere through vertical forcing only; alternatively, existing plate motions can induce new subduction zones by horizontally forcing the oceanic lithosphere to rupture and sink into the asthenosphere. Both models can eventually yield self-sustaining subduction zones, as the oceanic crust is metamorphosed at great depth and becomes denser than the surrounding mantle rocks. The compilation of subduction zone initiation events back to 100 Ma suggests horizontally-forced subduction zone initiation for most modern subduction zones, which is supported by results from numerical models and geologic studies. Some analogue modeling shows, however, the possibility of spontaneous subduction from inherent density differences between two plates at specific locations like passive margins and along transform faults. There is evidence this has taken place in the Izu-Bonin-Mariana subduction system. Earlier in Earth's history, subduction is likely to have initiated without horizontal forcing due to the lack of relative plate motion, though a proposal by A. Yin suggests that meteorite impacts may have contributed to subduction initiation on early Earth.
Though the idea of subduction initiation at passive margins is popular, there is no modern day example for this type of subduction nucleation. This is likely due to the strength of the oceanic or transitional crust at the continental passive margins, suggesting that if the crust did not break in its first 20 million years of life, it is unlikely to break in the future under normal sedimentation loads. Only with additional weaking of the crust, through hotspot magmatism or extensional rifting, would the crust be able to break from its continent and begin subduction.
End of subduction
Subduction can continue as long as the oceanic lithosphere moves into the subduction zone. However, the arrival of buoyant continental lithosphere at a subduction zone can result in increased coupling at the trench and cause plate boundary reorganization. The arrival of continental crust results in continental collision or terrane accretion that may disrupt subduction. Continental crust can subduct to depths of where it can reach a point of no return. Sections of crustal or intraoceanic arc crust greater than in thickness or oceanic plateau greater than in thickness can disrupt subduction. However, island arcs subducted end-on may cause only local disruption, while an arc arriving parallel to the zone can shut it down. This has happened with the Ontong Java Plateau and the Vitiaz Trench.
Characteristics and Effects
Metamorphism
Subduction zones host a unique variety of rock types created by the high-pressure, low-temperature conditions a subducting slab encounters during its descent. The metamorphic conditions the slab passes through in this process create and destroy water bearing (hydrous) mineral phases, releasing water into the mantle. This water lowers the melting point of mantle rock, initiating melting. Understanding the timing and conditions in which these dehydration reactions occur is key to interpreting mantle melting, volcanic arc magmatism, and the formation of continental crust.
A metamorphic facies is characterized by a stable mineral assemblage specific to a pressure-temperature range and specific starting material. Subduction zone metamorphism is characterized by a low temperature, high-ultrahigh pressure metamorphic path through the zeolite, prehnite-pumpellyite, blueschist, and eclogite facies stability zones of subducted oceanic crust. Zeolite and prehnite-pumpellyite facies assemblages may or may not be present, thus the onset of metamorphism may only be marked by blueschist facies conditions. Subducting slabs are composed of basaltic crust topped with pelagic sediments; however, the pelagic sediments may be accreted onto the forearc-hanging wall and not subducted. Most metamorphic phase transitions that occur within the subducting slab are prompted by the dehydration of hydrous mineral phases. The breakdown of hydrous mineral phases typically occurs at depths greater than 10 km. Each of these metamorphic facies is marked by the presence of a specific stable mineral assemblage, recording the metamorphic conditions undergone but the subducting slab. Transitions between facies cause hydrous minerals to dehydrate at certain pressure-temperature conditions and can therefore be tracked to melting events in the mantle beneath a volcanic arc.
Arc magmatism
Two kinds of arcs are generally observed on Earth: island arcs that form on the oceanic lithosphere (for example, the Mariana and the Tonga island arcs), and continental arcs such as the Cascade Volcanic Arc, that form along the coast of continents. Island arcs (intraoceanic or primitive arcs) are produced by the subduction of oceanic lithosphere beneath another oceanic lithosphere (ocean-ocean subduction) while continental arcs (Andean arcs) form during the subduction of oceanic lithosphere beneath a continental lithosphere (ocean-continent subduction). An example of a volcanic arc having both island and continental arc sections is found behind the Aleutian Trench subduction zone in Alaska.
Volcanoes that occur above subduction zones, such as Mount St. Helens, Mount Etna, and Mount Fuji, lie approximately one hundred kilometers from the trench in arcuate chains called volcanic arcs. Plutons, like Half Dome in Yosemite National Park, generally form 10–50 km below the volcanoes within the volcanic arcs and are only visible on the surface once the volcanoes have weathered away. The volcanism and plutonism occur as a consequence of the subducting oceanic slab dehydrating as it reaches higher pressures and temperatures. Once the oceanic slab reaches about 100 km in depth, hydrous minerals become unstable and release fluids into the asthenosphere. The fluids act as a flux for the rock within the asthenosphere and cause it to partially melt. The partially melted material is more buoyant and as a result will rise into the lithosphere, where it forms large magma chambers called diapirs. Some of the magma will make it to the surface of the crust where it will form volcanoes and, if eruptive on earth's surface, will produce andesitic lava. Magma that remains in the lithosphere long enough will cool and form plutonic rocks such as diorite, granodiorite, and sometimes granite.
The arc magmatism occurs one hundred to two hundred kilometers from the trench and approximately one hundred kilometers above the subducting slab. Arcs produce about 10% of the total volume of magma produced each year on Earth (approximately 0.75 cubic kilometers), much less than the volume produced at mid-ocean ridges, but they have formed most continental crust. Arc volcanism has the greatest impact on humans because many arc volcanoes lie above sea level and erupt violently. Aerosols injected into the stratosphere during violent eruptions can cause rapid cooling of Earth's climate and affect air travel.
Arc-magmatism plays a role in Earth's Carbon cycle by releasing subducted carbon through volcanic processes. Older theory states that the carbon from the subducting plate is made available in overlying magmatic systems via decarbonation, where CO is released through silicate-carbonate metamorphism. However, evidence from thermodynamic modeling has shown that the pressures and temperatures necessary for this type of metamorphism are much higher than what is observed in most subduction zones. Frezzoti et al. (2011) propose a different mechanism for carbon transport into the overriding plate via dissolution (release of carbon from carbon-bearing minerals into an aqueous solution) instead of decarbonation. Their evidence comes from the close examination of mineral and fluid inclusions in low-temperature (<600 °C) diamonds and garnets found in an eclogite facies in the Alps. The chemistry of the inclusions supports the existence of a carbon-rich fluid in that environment, and additional chemical measurements of lower pressure and temperature facies in the same tectonic complex support a model for carbon dissolution (rather than decarbonation) as a means of carbon transport.
Earthquakes and tsunamis
Elastic strain caused by plate convergence in subduction zones produces at least three types of earthquakes. These are deep earthquakes, megathrust earthquakes, and outer rise earthquakes. Deep earthquakes happen within the crust, megathrust earthquakes on the subduction interface near the trench, and outer rise earthquakes on the subducting lower plate as it bends near the trench.
Anomalously deep events are a characteristic of subduction zones, which produce the deepest quakes on the planet. Earthquakes are generally restricted to the shallow, brittle parts of the crust, generally at depths of less than twenty kilometers. However, in subduction zones quakes occur at depths as great as . These quakes define inclined zones of seismicity known as Wadati–Benioff zones which trace the descending slab.
Nine of the ten largest earthquakes of the last 100 years were subduction zone megathrust earthquakes. These included the 1960 Great Chilean earthquake which at M 9.5 was the largest earthquake ever recorded, the 2004 Indian Ocean earthquake and tsunami, and the 2011 Tōhoku earthquake and tsunami. The subduction of cold oceanic lithosphere into the mantle depresses the local geothermal gradient and causes a larger portion of Earth's crust to deform in a more brittle fashion than it would in a normal geothermal gradient setting. Because earthquakes can occur only when a rock is deforming in a brittle fashion, subduction zones can cause large earthquakes. If such a quake causes rapid deformation of the sea floor, there is potential for tsunamis. The largest tsunami ever recorded happened due to a mega-thrust earthquake on December 26, 2004. The earthquake was caused by subduction of the Indo-Australian plate under the Euro-Asian Plate, but the tsunami spread over most of the planet and devastated the areas around the Indian Ocean. Small tremors which cause small, nondamaging tsunamis, also occur frequently.
A study published in 2016 suggested a new parameter to determine a subduction zone's ability to generate mega-earthquakes. By examining subduction zone geometry and comparing the degree of lower plate curvature of the subducting plate in great historical earthquakes such as the 2004 Sumatra-Andaman and the 2011 Tōhoku earthquake, it was determined that the magnitude of earthquakes in subduction zones is inversely proportional to the angle of subduction near the trench, meaning that "the flatter the contact between the two plates, the more likely it is that mega-earthquakes will occur".
Outer rise earthquakes on the lower plate occur when normal faults oceanward of the subduction zone are activated by flexure of the plate as it bends into the subduction zone. The 2009 Samoa earthquake is an example of this type of event. Displacement of the sea floor caused by this event generated a six-meter tsunami in nearby Samoa.
Seismic tomography has helped detect subducted lithospheric slabs deep in the mantle where no earthquakes occur. About one hundred slabs have been described in terms of depth and their timing and location of subduction. The great seismic discontinuities in the mantle, at depth and , are disrupted by the descent of cold slabs in deep subduction zones. Some subducted slabs seem to have difficulty penetrating the major discontinuity that marks the boundary between the upper mantle and lower mantle at a depth of about 670 kilometers. Other subducted oceanic plates have sunk to the core–mantle boundary at 2890 km depth. Generally, slabs decelerate during their descent into the mantle, from typically several cm/yr (up to ~10 cm/yr in some cases) at the subduction zone and in the uppermost mantle, to ~1 cm/yr in the lower mantle. This leads to either folding or stacking of slabs at those depths, visible as thickened slabs in seismic tomography. Below ~1700 km, there might be a limited acceleration of slabs due to lower viscosity as a result of inferred mineral phase changes until they approach and finally stall at the core–mantle boundary. Here the slabs are heated up by the ambient heat and are not detected anymore ~300 Myr after subduction.
Orogeny
Orogeny is the process of mountain building. Subducting plates can lead to orogeny by bringing oceanic islands, oceanic plateaus, sediments and passive continental margins to convergent margins. The material often does not subduct with the rest of the plate but instead is accreted to (scraped off) the continent, resulting in exotic terranes. The collision of this oceanic material causes crustal thickening and mountain-building. The accreted material is often referred to as an accretionary wedge or prism. These accretionary wedges can be associated with ophiolites (uplifted ocean crust consisting of sediments, pillow basalts, sheeted dykes, gabbro, and peridotite).
Subduction may also cause orogeny without bringing in oceanic material that accretes to the overriding continent. When the lower plate subducts at a shallow angle underneath a continent (something called "flat-slab subduction"), the subducting plate may have enough traction on the bottom of the continental plate to cause the upper plate to contract by folding, faulting, crustal thickening, and mountain building. Flat-slab subduction causes mountain building and volcanism moving into the continent, away from the trench, and has been described in western North America (i.e. Laramide orogeny, and currently in Alaska, South America, and East Asia.
The processes described above allow subduction to continue while mountain building happens concurrently, which is in contrast to continent-continent collision orogeny, which often leads to the termination of subduction.
Subduction of continental lithosphere
Continents are pulled into subduction zones by the sinking oceanic plate they are attached to. Where continents are attached to oceanic plates with no subduction, there is a deep basin that accumulates thick suites of sedimentary and volcanic rocks known as a passive margin. Some passive margins have up to 10 km of sedimentary and volcanic rocks covering the continental crust. As a passive margin is pulled into a subduction zone by the attached and negatively buoyant oceanic lithosphere, the sedimentary and volcanic cover is mostly scraped off to form an orogenic wedge. An orogenic wedge is larger than most accretionary wedges due to the volume of material there is to accrete. The continental basement rocks beneath the weak cover suites are strong and mostly cold, and can be underlain by a >200 km thick layer of dense mantle. After shedding the low density cover units, the continental plate, especially if it is old, goes down the subduction zone. As this happens, metamorphic reactions increase the density of the continental crustal rocks, which leads to less buoyancy.
One study of the active Banda arc-continent collision claims that by unstacking the layers of rock that once covered the continental basement, but are now thrust over one another in the orogenic wedge, and measuring how long they are, can provide a minimum estimate of how far the continent has subducted. The results show at least a minimum of 229 kilometers of subduction of the northern Australian continental plate. Another example may be the continued northward motion of India, which is subducting beneath Asia. The collision between the two continents initiated around 50 my ago, but is still active.
Intra-oceanic: ocean/ocean plate subduction
Oceanic-Oceanic plate subduction zones comprise roughly 40% of all subduction zone margins on the planet. The ocean-ocean plate relationship can lead to subduction zones between oceanic and continental plates, therefore highlighting how important it is to understand this subduction setting. Although it is not fully understood what causes the initiation of subduction of an oceanic plate under another oceanic plate, there are three main models put forth by Baitsch-Ghirardello et al. that explain the different regimes present in this setting.
The models are as follows:
retreating subduction: caused by weak coupling between the lower and upper plate which leads to the opening of a back arc basin and the subduction zone being moved by slab rollback.
stable subduction: caused by intermediate coupling between the lower and upper plate. The subduction zone generally stays in the same place and the subduction plate subducts at a consistent angle.
advancing subduction: caused by strong coupling between the upper and lower plate. The subducting sediments thicken causing partially molten plumes to be on top of subducting plate.
Arc-continent collision and global climate
In their 2019 study, Macdonald et al. proposed that arc-continent collision zones and the subsequent obduction of oceanic lithosphere was at least partially responsible for controlling global climate. Their model relies on arc-continent collision in tropical zones, where exposed ophiolites composed mainly of mafic material increase "global weatherability" and result in the storage of carbon through silicate weathering processes. This storage represents a carbon sink, removing carbon from the atmosphere and resulting in global cooling. Their study correlates several Phanerozoic ophiolite complexes, including active arc-continent subduction, with known global cooling and glaciation periods. This study does not discuss Milankovitch cycles as a driver of global climate cyclicity.
Beginnings of subduction on Earth
Modern-style subduction is characterized by low geothermal gradients and the associated formation of high-pressure low-temperature rocks such as eclogite and blueschist. Likewise, rock assemblages called ophiolites, associated with modern-style subduction, also indicate such conditions. Eclogite xenoliths found in the North China Craton provide evidence that modern-style subduction occurred at least as early as 1.8 Ga ago in the Paleoproterozoic Era. The eclogite itself was produced by oceanic subduction during the assembly of supercontinents at about 1.9–2.0 Ga.
Blueschist is a rock typical for present-day subduction settings. The absence of blueschist older than Neoproterozoic reflects more magnesium-rich compositions of Earth's oceanic crust during that period. These more magnesium-rich rocks metamorphose into greenschist at conditions when modern oceanic crust rocks metamorphose into blueschist. The ancient magnesium-rich rocks mean that Earth's mantle was once hotter, but not that subduction conditions were hotter. Previously, the lack of pre-Neoproterozoic blueschist was thought to indicate a different type of subduction. Both lines of evidence refute previous conceptions of modern-style subduction having been initiated in the Neoproterozoic Era 1.0 Ga ago.
History of investigation
Harry Hammond Hess, who during World War II served in the United States Navy Reserve and became fascinated in the ocean floor, studied the Mid-Atlantic Ridge and proposed that hot molten rock was added to the crust at the ridge and expanded the seafloor outward. This theory was to become known as seafloor spreading. Since the Earth's circumference has not changed over geologic time, Hess concluded that older seafloor has to be consumed somewhere else, and suggested that this process takes place at oceanic trenches, where the crust would be melted and recycled into the Earth's mantle.
In 1964, George Plafker researched the Good Friday earthquake in Alaska. He concluded that the cause of the earthquake was a megathrust reaction in the Aleutian Trench, a result of the Alaskan continental crust overlapping the Pacific oceanic crust. This meant that the Pacific crust was being forced downward, or subducted, beneath the Alaskan crust. The concept of subduction would play a role in the development of the plate tectonics theory.
First geologic attestations of the "subduct" words date to 1970, In ordinary English to subduct, or to subduce (from Latin subducere, "to lead away") are transitive verbs requiring a subject to perform an action on an object not itself, here the lower plate, which has then been subducted ("removed"). The geological term is "consumed", which happens the geological moment the lower plate slips under, even though it may persist for some time until its remelting and dissipation. In this conceptual model, plate is continually being used up. The identity of the subject, the consumer, or agent of consumption, is left unstated. Some sources accept this subject-object construct.
Geology makes to subduct into an intransitive verb and a reflexive verb. The lower plate itself is the subject. It subducts, in the sense of retreat, or removes itself, and while doing so, is the "subducting plate". Moreover, the word slab is specifically attached to the "subducting plate", even though in English the upper plate is just as much of a slab. The upper plate is left hanging, so to speak. To express it geology must switch to a different verb, typically to override. The upper plate, the subject, performs the action of overriding the object, the lower plate, which is overridden.
Importance
Subduction zones are important for several reasons:
Subduction zone physics: Sinking of the oceanic lithosphere (sediments, crust, mantle), by the contrast of density between the cold and old lithosphere and the hot asthenospheric mantle wedge, is the strongest force (but not the only one) needed to drive plate motion and is the dominant mode of mantle convection.
Subduction zone chemistry: The subducted sediments and crust dehydrate and release water-rich (aqueous) fluids into the overlying mantle, causing mantle melting and fractionation of elements between the surface and deep mantle reservoirs, producing island arcs and continental crust. Hot fluids in subduction zones also alter the mineral compositions of the subducting sediments and potentially the habitability of the sediments for microorganisms.
Subduction zones drag down subducted oceanic sediments, oceanic crust, and mantle lithosphere that interact with the hot asthenospheric mantle from the over-riding plate to produce calc-alkaline series melts, ore deposits, and continental crust.
Subduction zones pose significant threats to lives, property, economic vitality, cultural and natural resources, and quality of life. The tremendous magnitudes of earthquakes and volcanic eruptions can also have knock-on effects with global impact.
Subduction zones have also been considered as possible disposal sites for nuclear waste in which the action of subduction itself would carry the material into the planetary mantle, safely away from any possible influence on humanity or the surface environment. However, that method of disposal is currently banned by international agreement. Furthermore, plate subduction zones are associated with very large megathrust earthquakes, making the effects of using any specific site for disposal unpredictable and possibly adverse to the safety of long-term disposal.
| Physical sciences | Tectonics | Earth science |
145865 | https://en.wikipedia.org/wiki/Parts-per%20notation | Parts-per notation | In science and engineering, the parts-per notation is a set of pseudo-units to describe small values of miscellaneous dimensionless quantities, e.g. mole fraction or mass fraction. Since these fractions are quantity-per-quantity measures, they are pure numbers with no associated units of measurement. Commonly used are parts-per-million (ppm, ), parts-per-billion (ppb, ), parts-per-trillion (ppt, ) and parts-per-quadrillion (ppq, ). This notation is not part of the International System of Units (SI) system and its meaning is ambiguous.
Applications
Parts-per notation is often used describing dilute solutions in chemistry, for instance, the relative abundance of dissolved minerals or pollutants in water. The quantity "1 ppm" can be used for a mass fraction if a water-borne pollutant is present at one-millionth of a gram per gram of sample solution. When working with aqueous solutions, it is common to assume that the density of water is 1.00 g/mL. Therefore, it is common to equate 1 kilogram of water with 1 L of water. Consequently, 1 ppm corresponds to 1 mg/L and 1 ppb corresponds to 1 μg/L.
Similarly, parts-per notation is used also in physics and engineering to express the value of various proportional phenomena. For instance, a special metal alloy might expand 1.2 micrometers per meter of length for every degree Celsius and this would be expressed as Parts-per notation is also employed to denote the change, stability, or uncertainty in measurements. For instance, the accuracy of land-survey distance measurements when using a laser rangefinder might be 1 millimeter per kilometer of distance; this could be expressed as "Accuracy = 1 ppm."
Parts-per notations are all dimensionless quantities: in mathematical expressions, the units of measurement always cancel. In fractions like "2 nanometers per meter" so the quotients are pure-number coefficients with positive values less than or equal to 1. When parts-per notations, including the percent symbol (%), are used in regular prose (as opposed to mathematical expressions), they are still pure-number dimensionless quantities. However, they generally take the literal "parts per" meaning of a comparative ratio (e.g. "2 ppb" would generally be interpreted as "two parts in a billion parts").
Parts-per notations may be expressed in terms of any unit of the same measure. For instance, the expansion coefficient of some brass alloy, may be expressed as 18.7 (μm/m)/°C, or as 18.7 (μ in/in)/°C; the numeric value representing a relative proportion does not change with the adoption of a different unit of length.
Similarly, a metering pump that injects a trace chemical into the main process line at the proportional flow rate is doing so at a rate that may be expressed in a variety of volumetric units, including 125 cm3/m3, etc.
In nuclear magnetic resonance spectroscopy (NMR), chemical shift is usually expressed in ppm. It represents the difference of a measured frequency in parts per million from the reference frequency. The reference frequency depends on the instrument's magnetic field and the element being measured. It is usually expressed in MHz. Typical chemical shifts are rarely more than a few hundred Hz from the reference frequency, so chemical shifts are conveniently expressed in ppm (Hz/MHz). Parts-per notation gives a dimensionless quantity that does not depend on the instrument's field strength.
Parts-per expressions
One part per hundred is generally represented by the percent sign (%) and denotes one part per 100 () parts, and a value of . This is equivalent to about fourteen minutes out of one day.
One part per thousand should generally be spelled out in full and not as "ppt" (which is usually understood to represent "parts per trillion"). It may also be denoted by the permille sign (‰). Note however, that specific disciplines such as oceanography, as well as educational exercises, do use the "ppt" abbreviation. "One part per thousand" denotes one part per 1,000 () parts, and a value of . This is equivalent to about ninety seconds out of one day.
One part per ten thousand is denoted by the permyriad sign (‱). Although rarely used in science (ppm is typically used instead), one permyriad has an unambiguous value of one part per 10,000 () parts, and a value of . This is equivalent to about nine seconds out of one day. In contrast, in finance, the basis point is typically used to denote changes in or differences between percentage interest rates (although it can also be used in other cases where it is desirable to express quantities in hundredths of a percent). For instance, a change in an interest rate from 5.15% per annum to 5.35% per annum could be denoted as a change of 20 basis points (per annum). As with interest rates, the words "per annum" (or "per year") are often omitted. In that case, the basis point is a quantity with a dimension of (time−1).
One part per hundred thousand, per cent mille (pcm) or milli-percent denotes one part per 100,000 () parts, and a value of . It is commonly used in epidemiology for mortality, crime and disease prevalence rates, and nuclear reactor engineering as a unit of reactivity. In time measurement it is equivalent to about 5 minutes out of a year; in distance measurement, it is equivalent to 1 cm of error per km of distance traversed.
One part per million (ppm) denotes one part per 1,000,000 () parts, and a value of . It is equivalent to about 32 seconds out of a year or 1 mm of error per km of distance traversed. In mining, it is also equivalent to one gram per metric ton, expressed as g/t.
One part per billion (ppb) denotes one part per 1,000,000,000 () parts, and a value of . This is equivalent to about three seconds out of a century.
One part per trillion (ppt) denotes one part per 1,000,000,000,000 () parts, and a value of . This is equivalent to about thirty seconds out of every million years.
One part per quadrillion (ppq) denotes one part per 1,000,000,000,000,000 () parts, and a value of . This is equivalent to about two and a half minutes out of the age of the Earth (4.5 billion years). Although relatively uncommon in analytical chemistry, measurements at the ppq level are sometimes performed.
Criticism
Although the International Bureau of Weights and Measures (an international standards organization known also by its French-language initials BIPM) recognizes the use of parts-per notation, it is not formally part of the International System of Units (SI). Note that although "percent" (%) is not formally part of the SI, both the BIPM and the International Organization for Standardization (ISO) take the position that "in mathematical expressions, the internationally recognized symbol % (percent) may be used with the SI to represent the number 0.01" for dimensionless quantities. According to IUPAP, "a continued source of annoyance to unit purists has been the continued use of percent, ppm, ppb, and ppt". Although SI-compliant expressions should be used as an alternative, the parts-per notation remains nevertheless widely used in technical disciplines. The main problems with the parts-per notation are set out below.
Long and short scales
Because the named numbers starting with a "billion" have different values in different countries, the BIPM suggests avoiding the use of "ppb" and "ppt" to prevent misunderstanding. The U.S. National Institute of Standards and Technology (NIST) takes the stringent position, stating that "the language-dependent terms [...] are not acceptable for use with the SI to express the values of quantities".
Thousand vs. trillion
Although "ppt" usually means "parts per trillion", it occasionally means "parts per thousand". Unless the meaning of "ppt" is defined explicitly, it has to be determined from the context.
Mass fraction vs. mole fraction vs. volume fraction
Another problem of the parts-per notation is that it may refer to mass fraction, mole fraction or volume fraction. Since it is usually not stated which quantity is used, it is better to write the units out, such as kg/kg, mol/mol or m3/m3, even though they are all dimensionless. The difference is quite significant when dealing with gases, and it is very important to specify which quantity is being used. For example, the conversion factor between a mass fraction of 1 ppb and a mole fraction of 1 ppb is about 4.7 for the greenhouse gas CFC-11 in air (Molar mass of CFC-11 / Mean molar mass of air = 137.368 / 28.97 = 4.74). For volume fraction, the suffix "V" or "v" is sometimes appended to the parts-per notation (e.g. ppmV, ppbv, pptv). However, ppbv and pptv are usually used to mean mole fractions"volume fraction" would literally mean what volume of a pure substance is included in a given volume of a mixture, and this is rarely used except in the case of alcohol by volume.
To distinguish the mass fraction from volume fraction or mole fraction, the letter "w" (standing for "weight") is sometimes added to the abbreviation (e.g. ppmw, ppbw).
The usage of the parts-per notation is generally quite fixed within each specific branch of science, but often in a way that is inconsistent with its usage in other branches, leading some researchers to assume that their own usage (mass/mass, mol/mol, volume/volume, mass/volume, or others) is correct and that other usages are incorrect. This assumption sometimes leads them to not specify the details of their own usage in their publications, and others may therefore misinterpret their results. For example, electrochemists often use volume/volume, while chemical engineers may use mass/mass as well as volume/volume, while chemists, the field of occupational safety and the field of permissible exposure limit (e.g. permitted gas exposure limit in air) may use mass/volume. Unfortunatelly, many academic publications of otherwise excellent level fail to specify their use of the parts-per notation, which irritates some readers, especially those who are not experts in the particular fields in those publications, because parts-per-notation, without specifying what it stands for, can mean anything.
SI-compliant expressions
SI-compliant units that can be used as alternatives are shown in the chart below. Expressions that the BIPM explicitly does not recognize as being suitable for denoting dimensionless quantities with the SI are marked with !.
Note that the notations in the "SI units" column above are for the most part dimensionless quantities; that is, the units of measurement factor out in expressions like "1 nm/m" (1 nm/m =1 × 10−9) so the ratios are pure-number coefficients with values less than 1.
Uno (proposed dimensionless unit)
Because of the cumbersome nature of expressing certain dimensionless quantities per SI guidelines, the International Union of Pure and Applied Physics (IUPAP) in 1999 proposed the adoption of the special name "uno" (symbol: U) to represent the number 1 in dimensionless quantities. In 2004, a report to the International Committee for Weights and Measures (CIPM) stated that the response to the proposal of the uno "had been almost entirely negative", and the principal proponent "recommended dropping the idea". To date, the uno has not been adopted by any standards organization.
| Physical sciences | Concentration | Basics and measurement |
146000 | https://en.wikipedia.org/wiki/Amygdala | Amygdala | The amygdala (; : amygdalae or amygdalas; also ; Latin from Greek, , , 'almond', 'tonsil') is a paired nuclear complex present in the cerebral hemispheres of vertebrates. It is considered part of the limbic system. In primates, it is located medially within the temporal lobes. It consists of many nuclei, each made up of further subnuclei. The subdivision most commonly made is into the basolateral, central, cortical, and medial nuclei together with the intercalated cell clusters. The amygdala has a primary role in the processing of memory, decision-making, and emotional responses (including fear, anxiety, and aggression). The amygdala was first identified and named by Karl Friedrich Burdach in 1822.
Structure
Thirteen nuclei have been identified, each with their own subdivisions and distinct connections to the rest of the brain. The chief nuclei are the basolateral complex, the central nucleus, the cortical nucleus, the medial nucleus, and the intercalated cell clusters.
The basolateral complex can be further subdivided into the lateral, the basal, and the accessory basal nuclei. It has extensive connections with higher-order cortical areas in the prefrontal, temporal, insular cortices, and the hippocampus.
The basolateral complex is surrounded the intercalated cell net that is inhibitory and projects to a broad variety of areas in the basal forebrain, hypothalamus, and the amygdala.
The cortical and medial nuclei connect with the olfactory system and hypothalamus.
The central nucleus has extensive projections to the brainstem.
According to Larry Swanson and Gorica Petrovich, in an article titled, What is the amygdala? "The amygdala is neither a structural nor a functional unit".
Hemispheric specializations
In one study, electrical stimulations of the right amygdala induced negative emotions, especially fear and sadness. In contrast, stimulation of the left amygdala was able to induce either pleasant (happiness) or unpleasant (fear, anxiety, sadness) emotions. Other evidence suggests that the left amygdala plays a role in the brain's reward system.
Each side holds a specific function in how we perceive and process emotion. The right and left portions of the amygdala have independent memory systems, but work together to store, encode, and interpret emotion.
The right hemisphere of the amygdala is associated with negative emotion. It plays a role in the expression of fear and in the processing of fear-inducing stimuli. Fear conditioning, which occurs when a neutral stimulus acquires aversive properties, occurs within the right hemisphere. When an individual is presented with a conditioned, aversive stimulus, it is processed within the right amygdala, producing an unpleasant or fearful response. This emotional response conditions the individual to avoid fear-inducing stimuli and more importantly, to assess threats in the environment.
The right hemisphere is also linked to declarative memory, which consists of facts and information from previously experienced events and must be consciously recalled. It also plays a significant role in the retention of episodic memory. Episodic memory consists of the autobiographical aspects of memory, permitting recall of emotional and sensory experience of an event. This type of memory does not require conscious recall. The right amygdala plays a role in the association of time and places with emotional properties.
Development and sex distinction
The amygdala is one of the best-understood brain regions with regard to differences between the sexes. The amygdala is larger in males than females, in children aged 7 to 11, adult humans, and adult rats.
There is considerable growth within the first few years of structural development in both male and female amygdalae. Within this early period, female limbic structures grow at a more rapid pace than the male ones. Amongst female subjects, the amygdala reaches its full growth potential approximately 1.5 years before the peak of male development. The structural development of the male amygdala occurs over a longer period than in women. Because of the early development of female amygdalae, they reach their growth potential sooner than males, whose amygdalae continue to develop. The larger relative size of the male amygdala may be attributed to this extended developmental period.
Hormonal factors may contribute to these sex-specific developmental differences. The amygdala is rich in androgen receptors – nuclear receptors that bind to testosterone. Androgen receptors play a role in the DNA binding that regulates gene expression. Though testosterone is present within the female hormonal systems, women have lower levels of testosterone than men. The abundance of testosterone in the male hormonal system may contribute to development. In addition, the grey matter volume on the amygdala is predicted by testosterone levels, which may also contribute to the increased mass of the male amygdala.
There are observable developmental differences between the right and left amygdala. The left amygdala reaches its developmental peak approximately 1.5–2 years prior to the right amygdala. Despite the early growth of the left amygdala, the right increases in volume for a longer period of time. The right amygdala is associated with response to fearful stimuli as well as face recognition. It is inferred that the early development of the left amygdala functions to provide infants the ability to detect danger. In childhood, the amygdala is found to react differently to same-sex versus opposite-sex individuals. This reactivity decreases until a person enters adolescence, where it increases dramatically at puberty.
Other functional and structural differences between male and female amygdalae have been observed. Subjects' amygdala activation was observed when watching a horror film and subliminal stimuli. The results of the study showed a different lateralization of the amygdala in men and women. Enhanced memory for the film was related to enhanced activity of the left, but not the right, amygdala in women, whereas it was related to enhanced activity of the right, but not the left, amygdala in men. Similarly, a study of decision-making ability in patients with unilateral amygdala damage suggested that men with right (but not left) amygdala damage were more likely to be impaired in decision-making ability, while women with left (but not right) amygdala damage were more likely to be impaired in decision-making ability. One study found evidence that on average, women tend to retain stronger memories for emotional events than men.
Function
Connections
A simple view of the information processing through the amygdala follows as: the amygdala sends projections to the hypothalamus, septal nuclei and BNST (via the amygdalofugal tract), the dorsomedial thalamus (via the amygdalothalamic tract), the nuclei of the trigeminal nerve and the facial nerve, the ventral tegmental area, the locus coeruleus, and the laterodorsal tegmental nucleus. The basolateral amygdala projects to the nucleus accumbens, including the medial shell.
The medial nucleus is involved in the sense of smell and pheromone-processing. It receives input from the olfactory bulb and olfactory cortex. The lateral amygdalae, which send impulses to the rest of the basolateral complexes and to the centromedial nuclei, receive input from the sensory systems. The centromedial nuclei are the main outputs for the basolateral complexes, and are involved in emotional arousal in rats and cats.
Variability in amygdala connectivity has been related to a variety of behaviors and outcomes such as fear recognition and social network size.
Emotional learning
In complex vertebrates, including humans, the amygdalae perform primary roles in the formation and storage of memories associated with emotional events. Research indicates that, during fear conditioning, sensory stimuli reach the basolateral complexes of the amygdalae, particularly the lateral nuclei, where they form associations with memories of the stimuli. The association between stimuli and the aversive events they predict may be mediated by long-term potentiation, a sustained enhancement of signaling between affected neurons. There have been studies that show that damage to the amygdala can interfere with memory that is strengthened by emotion. One study examined a patient with bilateral degeneration of the amygdala. He was told a violent story accompanied by matching pictures and was observed based on how much he could recall from the story. The patient had less recollection of the story than patients with functional amygdala, showing that the amygdala has a strong connection with emotional learning.
Emotional memories are thought to be stored in synapses throughout the brain. Fear memories, for example, are considered to be stored in the neuronal connections from the lateral nuclei to the central nucleus of the amygdalae and the bed nuclei of the stria terminalis (part of the extended amygdala). These connections are not the sole site of fear memories given that the nuclei of the amygdala receive and send information to other brain regions that are important for memory such as the hippocampus. Some sensory neurons project their axon terminals to the central nucleus. The central nuclei are involved in the genesis of many fear responses such as defensive behavior (freezing or escape responses), autonomic nervous system responses (changes in blood pressure and heart rate/tachycardia), neuroendocrine responses (stress-hormone release), etc. Damage to the amygdalae impairs both the acquisition and expression of Pavlovian fear conditioning, a form of classical conditioning of emotional responses. Accumulating evidence has suggested that multiple neuromodulators acting in the amygdala regulates the formation of emotional memories.
The amygdalae are also involved in appetitive (positive) conditioning. It seems that distinct neurons respond to positive and negative stimuli, but there is no clustering of these distinct neurons into clear anatomical nuclei. However, lesions of the central nucleus in the amygdala have been shown to reduce appetitive learning in rats. Lesions of the basolateral regions do not exhibit the same effect. Research like this indicates that different nuclei within the amygdala have different functions in appetitive conditioning.
Nevertheless, researchers found an example of appetitive emotional learning showing an important role for the basolateral amygdala: The naïve female mice are innately attracted to non-volatile pheromones contained in male-soiled bedding, but not by the male-derived volatiles, become attractive if associated with non-volatile attractive pheromones, which act as unconditioned stimulus in a case of Pavlovian associative learning. In the vomeronasal, olfactory and emotional systems, Fos (gene family) proteins show that non-volatile pheromones stimulate the vomeronasal system, whereas air-borne volatiles activate only the olfactory system. Thus, the acquired preference for male-derived volatiles reveals an olfactory-vomeronasal associative learning. Moreover, the reward system is differentially activated by the primary pheromones and secondarily attractive odorants. Exploring the primary attractive pheromone activates the basolateral amygdala and the shell of nucleus accumbens but neither the ventral tegmental area nor the orbitofrontal cortex. In contrast, exploring the secondarily attractive male-derived odorants involves activation of a circuit that includes the basolateral amygdala, prefrontal cortex and ventral tegmental area. Therefore, the basolateral amygdala stands out as the key center for vomeronasal-olfactory associative learning.
Social
Reward
Glutamatergic neurons in the basolateral amygdala send projections to the nucleus accumbens shell and core. Activation of these projections drive motivational salience. The ability of these projections to drive incentive salience is dependent upon dopamine receptor D1.
Memory modulation
The amygdala is also involved in the modulation of memory consolidation. Following any learning event, the long-term memory for the event is not formed instantaneously. Rather, information regarding the event is slowly assimilated into long-term (potentially lifelong) storage over time, possibly via long-term potentiation. Recent studies suggest that the amygdala regulates memory consolidation in other brain regions. Also, fear conditioning, a type of memory that is impaired following amygdala damage, is mediated in part by long-term potentiation.
During the consolidation period, the memory can be modulated. In particular, it appears that emotional arousal following the learning event influences the strength of the subsequent memory for that event. Greater emotional arousal following a learning event enhances a person's retention of that event. Experiments have shown that administration of stress hormones to mice immediately after they learn something enhances their retention when they are tested two days later.
The amygdala, especially the basolateral nuclei, are involved in mediating the effects of emotional arousal on the strength of the memory for the event, as shown by many laboratories including that of James McGaugh. These laboratories have trained animals on a variety of learning tasks and found that drugs injected into the amygdala after training affect the animals' subsequent retention of the task. These tasks include basic classical conditioning tasks such as inhibitory avoidance, where a rat learns to associate a mild footshock with a particular compartment of an apparatus, and more complex tasks such as spatial or cued water maze, where a rat learns to swim to a platform to escape the water. If a drug that activates the amygdalae is injected into the amygdalae, the animals had better memory for the training in the task. If a drug that inactivates the amygdalae is injected, the animals had impaired memory for the task.
In rats, DNA damage was found to increase in the amygdala immediately after exposure to stress. Stress was induced by 30 minutes of restraint or by forced swimming. By seven days after exposure to these stresses, increased DNA damage was no longer detectable in the amygdala, probably because of DNA repair.
Buddhist monks who do compassion meditation have been shown to modulate their amygdala, along with their temporoparietal junction and insula, during their practice. In an fMRI study, more intensive insula activity was found in expert meditators than in novices.
Amygdala activity at the time of encoding information correlates with retention for that information. However, this correlation depends on the relative "emotionalness" of the information. More emotionally arousing information increases amygdalar activity, and that activity correlates with retention. Amygdala neurons show various types of oscillation during emotional arousal, such as theta activity. These synchronized neuronal events could promote synaptic plasticity (which is involved in memory retention) by increasing interactions between neocortical storage sites and temporal lobe structures involved in declarative memory.
Research using Rorschach test blot 03 finds that the number of unique responses to this random figure links to larger sized amygdalae. The researchers note, "Since previous reports have indicated that unique responses were observed at higher frequency in the artistic population than in the nonartistic normal population, this positive correlation suggests that amygdalar enlargement in the normal population might be related to creative mental activity."
Neuropsychological correlates of amygdala activity
Early research on primates provided explanations as to the functions of the amygdala, as well as a basis for further research. As early as 1888, rhesus monkeys with a lesioned temporal cortex (including the amygdala) were observed to have significant social and emotional deficits. Heinrich Klüver and Paul Bucy later expanded upon this same observation by showing that large lesions to the anterior temporal lobe produced noticeable changes, including overreaction to all objects, hypoemotionality, loss of fear, hypersexuality, and hyperorality, a condition in which inappropriate objects are placed in the mouth. Some monkeys also displayed an inability to recognize familiar objects and would approach animate and inanimate objects indiscriminately, exhibiting a loss of fear towards the experimenters. This behavioral disorder was later named Klüver-Bucy syndrome accordingly, and later research proved it was specifically due to amygdala lesions. Monkey mothers who had amygdala damage showed a reduction in maternal behaviors towards their infants, often physically abusing or neglecting them. In 1981, researchers found that selective radio frequency lesions of the whole amygdala caused Klüver-Bucy syndrome.
With advances in neuroimaging technology such as MRI, neuroscientists have made significant findings concerning the amygdala in the human brain. A variety of data shows the amygdala has a substantial role in mental states, and is related to many psychological disorders. Some studies have shown children with anxiety disorders tend to have a smaller left amygdala. In the majority of the cases, there was an association between an increase in the size of the left amygdala with the use of SSRIs (antidepressant medication) or psychotherapy. The left amygdala has been linked to social anxiety disorder, obsessive and compulsive disorders, and posttraumatic stress disorder, as well as more broadly to separation and generalized anxiety disorder. In a 2003 study, subjects with borderline personality disorder showed significantly greater left amygdala activity than normal control subjects. Some borderline patients even had difficulties classifying neutral faces or saw them as threatening. Individuals with psychopathy show reduced autonomic responses to instructed fear cues than otherwise healthy individuals. In 2006, researchers observed hyperactivity in the amygdala when patients were shown threatening faces or confronted with frightening situations. Patients with severe social phobia showed a correlation with increased response in the amygdala. Similarly, depressed patients showed exaggerated left amygdala activity when interpreting emotions for all faces, and especially for fearful faces. This hyperactivity was normalized when patients were administered antidepressant medication. By contrast, the amygdala has been observed to respond differently in people with bipolar disorder. A 2003 study found that adult and adolescent bipolar patients tended to have considerably smaller amygdala volumes and somewhat smaller hippocampal volumes. Many studies have focused on the connections between the amygdala and autism.
Studies in 2004 and 2006 showed that normal subjects exposed to images of frightened faces or faces of people from another race will show increased activity of the amygdala, even if that exposure is subliminal. However, the amygdala is not necessary for the processing of fear-related stimuli, since persons in whom it is bilaterally damaged show rapid reactions to fearful faces, even in the absence of a functional amygdala.
Sexual orientation
Recent studies have suggested possible correlations between brain structure, including differences in hemispheric ratios and connection patterns in the amygdala, and sexual orientation. Homosexual men tend to exhibit more feminine patterns in the amygdala than heterosexual males do, just as homosexual women tend to show more masculine patterns in the amygdala than heterosexual females do. It was observed that amygdala connections were more widespread from the left amygdala in homosexual males, as is also found in heterosexual females. Amygdala connections were more widespread from the right amygdala in homosexual females, as in heterosexual males.
Social
Increased activity in the amygdala following compassion-oriented meditation may contribute to social connectedness. Similarly, the structural white matter connectivity to other brain regions is also associated with social network size.
Amygdala volume correlates positively with both the size (the number of contacts a person has) and the complexity (the number of different groups to which a person belongs) of social networks. Individuals with larger amygdalae had larger and more complex social networks. The amygdala is responsible for facial recognition and allows others to respond appropriately to different emotional expressions. They were also better able to make accurate social judgments about other persons' faces. The amygdala's role in the analysis of social situations stems specifically from its ability to identify and process changes in facial features. It does not, however, process the direction of the gaze of the person being perceived.
The amygdala is also thought to be a determinant of the level of a person's emotional intelligence. It is particularly hypothesized that larger amygdalae allow for greater emotional intelligence, enabling greater societal integration and cooperation with others.
The amygdala processes reactions to violations concerning personal space. These reactions are absent in persons in whom the amygdala is damaged bilaterally. Furthermore, the amygdala is found to be activated in fMRI when people observe that others are physically close to them, such as when a person being scanned knows that an experimenter is standing immediately next to the scanner, versus standing at a distance.
Aggression
Animal studies have shown that stimulating the amygdala appears to increase both sexual and aggressive behavior. Likewise, studies using brain lesions have shown that harm to the amygdala may produce the opposite effect. Thus, it appears that this part of the brain may play a role in the display and modulation of aggression.
Fear
There are cases of human patients with focal bilateral amygdala lesions due to the rare genetic condition Urbach-Wiethe disease. Such patients fail to exhibit fear-related behaviors, leading one, S.M., to be dubbed the "woman with no fear". This finding reinforces the conclusion that the amygdala "plays a pivotal role in triggering a state of fear".
Alcoholism and binge drinking
The amygdala appears to play a role in binge drinking, being damaged by repeated episodes of intoxication and withdrawal. Protein kinase C-epsilon in the amygdala is important for regulating behavioral responses to morphine, ethanol, and controlling anxiety-like behavior. The protein is involved in controlling the function of other proteins and plays a role in development of the ability to consume a large amount of ethanol. The duration of chronic alcohol consumption and abstinence may affect dynamic brain network adaptations. When excessive drinking occurs, the amygdala is affected through behavioral changes and reduces the brain's plasticity. Brain plasticity is how our brain grows and develops; it is also how our neurons can make connections with other neurons. This ultimately increases our neural pathways allowing us to increase our knowledge of the world around us. When our brain plasticity decreases, it makes it difficult for neurons to make connections to other neurons. Often when binge drinking, or alcoholism occurs, our amygdala is affected and leads to behavior damage. These behavioral damages can be lack of control, inability to conduct oneself in a mature manner, aggressive behavior, loss of conduct, anxiety, depression, personality disorders, excessive drug intake, bi-polar disorder, confusion, higher tolerance levels, irritability, and inappropriate sexual behaviors with others and self.
Anxiety
There may also be a link between the amygdala and anxiety. In particular, there is a higher prevalence of females that are affected by anxiety disorders. In an experiment, degu pups were removed from their mother but allowed to hear her call. In response, the males produced increased serotonin receptors in the amygdala but females lost them. This led to the males being less affected by the stressful situation.
The clusters of the amygdala are activated when an individual expresses feelings of fear or aggression. This occurs because the amygdala is the primary structure of the brain responsible for fight or flight response. Anxiety and panic attacks can occur when the amygdala senses environmental stressors that stimulate fight or flight response. The amygdala is directly associated with conditioned fear. Conditioned fear is the framework used to explain the behavior produced when an originally neutral stimulus is consistently paired with a stimulus that evokes fear. The amygdala represents a core fear system in the human body, which is involved in the expression of conditioned fear. Fear is measured by changes in autonomic activity including increased heart rate, increased blood pressure, as well as in simple reflexes such as flinching or blinking.
The central nucleus of the amygdala has direct correlations to the hypothalamus and brainstem – areas directly related to fear and anxiety. This connection is evident from studies of animals that have undergone amygdalae removal. Such studies suggest that animals lacking an amygdala have less fear expression and indulge in non-species-like behavior. Many projection areas of the amygdala are critically involved in specific signs that are used to measure fear and anxiety.
Mammals have very similar ways of processing and responding to danger. Scientists have observed similar areas in the brain – specifically in the amygdala – lighting up or becoming more active when a mammal is threatened or beginning to experience anxiety. Similar parts of the brain are activated when rodents and humans alike observe a dangerous situation, the amygdala playing a crucial role in this assessment. By observing the amygdalae's functions, it can determine why one rodent may be much more anxious than another. There is a direct relationship between the activation of the amygdala and the level of anxiety the subject feels.
Feelings of anxiety start with a catalyst – an environmental stimulus that provokes stress. This can include various smells, sights, and internal sensations that result in anxiety. The amygdala reacts to this stimuli by preparing to either stand and fight or to turn and run. This response is triggered by the release of adrenaline into the bloodstream. Consequently, blood sugar rises, becoming immediately available to the muscles for quick energy. Shaking may occur in an attempt to return blood to the rest of the body. Apart from initiation of stress, long-term changes in amygdala neurons may also increase anxiety after long-term or traumatic stress, led by the action of stress-related hormones within the amygdala. On the flip side, blocking the action of stress hormones in the amygdala reduces anxiety. A better understanding of the amygdala and its various functions may lead to a new way of treating clinical anxiety.
Posttraumatic stress disorder
There seems to be a connection with the amygdalae and how the brain processes posttraumatic stress disorder. Multiple studies have found that the amygdalae may be responsible for the emotional reactions of PTSD patients. One study in particular found that when PTSD patients are shown pictures of faces with fearful expressions, their amygdalae tended to have a higher activation than someone without PTSD.
Bipolar disorder
Amygdala dysfunction during face emotion processing is well-documented in bipolar disorder. Individuals with bipolar disorder showed greater amygdala activity (especially the amygdala/medial-prefrontal-cortex circuit).
Additional images
| Biology and health sciences | Nervous system | Biology |
146031 | https://en.wikipedia.org/wiki/Himalayan%20cat | Himalayan cat | The Himalayan (short for Himalayan Persian, or Colourpoint Persian as it is commonly referred to in Europe), is a breed or sub-breed of long-haired cat similar in type to the Persian, with the exception of its blue eyes and its point colouration, which were derived from crossing the Persian with the Siamese. Some registries may classify the Himalayan as a long-haired sub-breed of Siamese, or a colorpoint sub-breed of Persian. The World Cat Federation has merged them with the Colorpoint Shorthair and Javanese into a single breed, the Colorpoint.
History
Work to formally establish a breed with combined Persian and Siamese traits, explicitly for the cat fancy, began in the United States in the 1930s at Harvard University, under the term Siamese–Persian, and the results were published in the Journal of Heredity in 1936, but were not adopted as a recognized breed by any major fancier groups at the time. Brian Sterling-Webb independently developed the cross-breed over a period of ten years in the UK, and in 1955 it was recognized there as the Longhaired Colourpoint by the Governing Council of the Cat Fancy (GCCF).
California cat breeder Jean Mill took a series of graduate classes in genetics at the University of California, Davis. By 1948, she was one of three breeders independently crossing the Persian and Siamese to create the Himalayan cat.
Separate US-based breeding efforts had begun around 1950, and a breeder known to sources simply as Mrs. Goforth received breed recognition from the Cat Fanciers' Association (CFA) near the end of 1957 for the Himalayan. Early breeders were mostly interested in adding Siamese colouration to long-haired cats, and therefore reinforced the stock by outbreeding to Persians only to retain the Persian trait dominance. However, by the 1960s, some were re-introducing Siamese stock and producing less "Persian-style" cats, In the 1980s, a concerted effort to re-establish the breed along more formally Persian lines ultimately caused the breed to be merged into Persian as a variant in some registries (e.g. in 1984 by CFA), and a decline in the "old" or Siamese-like specimens.
Recognition
The Himalayan is considered a colour variant of the Persian and not a separate breed by the Cat Fanciers' Association and the GCCF. The Himalayan is considered a separate breed by the American Cat Fanciers Association and The International Cat Association.
Appearance
The Himalayan resembles the Persian in type, conformation, and coat length and texture. The Himalayan does not resemble the Siamese in type.
Body
The Himalayan is medium to large in size with a cobby body and low legs.
Head
The Himalayan's head is round and massive with a round face and a thick neck. The nose is snubbed and pushed in.
Ears
The ears of the Himalayan are small and round tipped and slightly pointed forward.
Eyes
The eyes are large and round and spread well apart. Pointed Himalayans have blue eyes, non-pointed Himalayans have copper eyes except for the silver and golden tabby which have green eyes.
Coat
The Himalayan has a long and thick coat all over the body including the tail and ear and toe tufts.
Coat colours
The Himalayan comes in most colours with prohibited colours being mink and sepia.
Health
Like the Persian, the Himalayan is a brachycephalic breed which predisposes it to health issues such as respiratory infections, epiphora, corneal abrasions, ulcers, and corneal sequestration. Himalayans are also suspecitible to polycystic kidney disease, a hereditary condition that results in cysts growing in the kidney. Himalayans have a higher incidence of feline asthma.
In a review of over 5,000 cases of urate urolithiasis the Himalayan was under-represented, with an odds ratio of 0.37. A study of cats presented to the University of Missouri-Columbia Veterinary Medical Teaching Hospital that underwent radiography found 4 Himalayans out of a population of 16 to have hip dysplasia, higher than the 6.6% average for all cats.
Himalayans are predisposed to dermatophytosis (ringworm).
The Himalayan is predisposed to urticaria pigmentosa, a type of benign mast cell disorder.
Idiopathic facial dermatitis, also known as facial dermatitis of the Persian and Himalayan cat is a type of dermatitis only observed in the Persian and Himalayan cat. It's characterised by greasy skin, debris adhering to the folds of the face and nose, ceruminous otitis externa, secondary bacterial folliculitis and Malassezia dermatitis, and pruritus. Onset is at 10 months to 6 years.
In popular culture
In the CBS television detective series "Tucker's Witch" (1982), a Himalayan cat named Dickens is the familiar to witch Amanda Tucker. Amanda Tucker has a telepathic link with Dickens, who provides her and her husband with clairvoyant clues to help them solve mysteries. Dickens is featured prominently in the show's opening and closing credits.
In the spoof film Date Movie (2006), Mr. Jinxers is a parody of his Meet the Parents counterpart.
In the movies Homeward Bound: The Incredible Journey (1993) and Homeward Bound II: Lost in San Francisco (1996), one of the main characters is a Himalayan cat named Sassy (voiced by Sally Field).
The main character of the anime/manga Prince of Tennis, Ryoma Echizen, owns a Himalayan cat named Karupin (or Kalpin in the English translation).
Martha Stewart owns three Himalayans, named after composers: Beethoven, Mozart and Bartók. The cats have been featured in her commercials for Kmart, on her television show, Martha Stewart Living, and in her magazine, such as the cover of the February 1999 issue.
A Himalayan named Luna The Fashion Kitty became a social media phenomenon in 2011 with a popular Facebook page, a website, and several media references.
A Himalayan-Persian named Colonel Meow became an Internet celebrity in 2012, and entered Guinness World Records 2014 as the cat with the longest fur.
Mr. Jinx (also known as Jinxy, or simply just Jinx) from the Meet the Parents trilogy is a seal-point peke-faced Himalayan with an all-black tail.
The "narrator" of David Michie's series of books that begins with "The Dalai Lama's Cat" is a Himalayan cat.
A community-created cosmetic item for the Medic and Spy classes in the 2007 computer game Team Fortress 2 is a Red Point Himalayan cat named "Harry".
Gallery
| Biology and health sciences | Cats | Animals |
146038 | https://en.wikipedia.org/wiki/Devon%20Rex | Devon Rex | The Devon Rex is a tall-eared, short-haired breed of cat that emerged in England during the late 1950s. The breed is known for its atypical appearance, with an oddly shaped head, large eyes, and the short and wavy coat.
History
Origin
Beryl Cox came across a novel curly-coated kitten in Buckfastleigh, Devon in 1960 whom she decided to name Kirlee. Originally, Cox believed the cat's gene to be related to the Cornish Rex which led to her breeding Kirlee with Cornish Rexes. However, when the queen gave birth the kittens all had straight coats, which led to the discovery that Kirlee had a different mutation than Cornish Rexes. Following this Cox began a breeding program to try and preserve Kirlee's unique mutation.
Appearance
The Devon Rex is a very short haired breed with a medium build and a unique head type which gives the breed a 'pixie-like' appearance. The head is short with a broad wedge and the brow curving to a flat skull. The eyes are large, set wide, and oval-shaped. Devon Rexes may have any eye colour. The ears are large and set wide apart with rounded tips. The coat is short dense and soft and curls inwards a little giving it a waved or rippled effect. Some areas of the body such as the neck may lack enough fur for the wave/ripple effect. The whiskers and eyebrows are crinkled and twisted. Devon Rexes may come in any colour.
Health
Hereditary myopathy is found in some Devon Rexes. It is caused by a genetic variant known as COLQ and it appears anytime from 3 to 23 weeks of age. Typically, there is a chance that the myopathy in the cats might stabilize, however, most of the time, the condition worsens and the Devon rex cats die from laryngospasm, after obstructing their larynx with food.
The Devon Rex was found to be predisposed to feline atopic dermatitis in a retrospective study of cases of the disease.
The Devon Rex is predisposed to congenital hypotrichosis and Malassezia dermatitis.
Gallery
| Biology and health sciences | Cats | Animals |
146048 | https://en.wikipedia.org/wiki/Birman | Birman | The Birman, also called the "Sacred Cat of Burma", is a domestic cat breed. The Birman is a long-haired, colour-pointed cat distinguished by a silky coat, deep blue eyes, and contrasting white "gloves" on each paw.
The breed name is derived from Birmanie, the French form of Burma. The breed was first recognised in 1925 in France.
History
No clear record of the breed's origin exists. They are most often claimed to have originated as the companions of temple priests in northern Burma in the Mount of Lugh. Many stories exist of how the cats first came to France, including pairs of cats being given as a reward for helping defend a temple, or being smuggled out of Burma and Sweden by a Vanderbilt. Another pair of Birmans (or a pregnant female called Poupée de Maldapour) were said to have been stolen and later imported to France by Thadde Haddisch. The first traces of historical Birmans go back to a Mme Leotardi in Nice, France.
Birmans were almost wiped out as a breed during World War II. Only two cats were alive in Europe at the end of the war, a pair named Orloff and Xenia de Kaabaa, both belonging to Baudoin-Crevoisier. The foundation of the breed in postwar France were offspring of this pair. They had to be heavily outcrossed with long-hair breeds such as Persian and Siamese to rebuild the Birman breed. By the early 1950s, pure Birman litters were again being produced. The restored breed was recognized in Britain in 1965 and by the CFA in 1966.
The first Birman cats were seal point. The blue point colour was introduced in 1959 using blue Persian lines. New colours were later added by English breeders including chocolate, red, and tabby/lynx points. Birmans have also been used in the development of new breeds such as the Ragdoll.
Breed recognition
The Birman breed was first recognized in France by the Cat Club de France in 1925, then in England by the Governing Council of the Cat Fancy (GCCF) in 1966 and in United States by the Cat Fanciers' Association (CFA) in 1967. It was also recognized by the Canadian Cat Association (CCA) and the International Cat Association (TICA) in 1979.
Appearance
Birmans have a medium-sized, rectangular body with a broad face and distinct Roman nose. Their ears are ideally as wide on the base as they are tall and should be set as much on top of the head as on the side. The eyes are rounded and should be a deep sapphire blue.
The Birman's fur is medium-long and should have a silky texture. Unlike a Persian or Himalayan, they have no undercoat, thus are much less prone to matting. Coat colour is always pointed, save for the contrasting pure white, symmetrical "gloves" on each paw that are the trademark of the breed. The white must involve all toes and in front must stop at the articulation or at the transition of toes to metacarpals. These gloves should extend noticeably further up the back of the leg (referred to as the "laces"), finishing with an inverted V extended half to three-fourths up the hock. Any other spot of white on the points is considered a serious fault. The base body colour is white to cream, with a wash of colour that corresponds to the points, but is much paler.
Recognized point colours are seal, chocolate, red and the corresponding dilute varieties: blue, lilac and cream. Tabby and tortie variations in seal, chocolate, blue or lilac are also allowed; other colours are in development.
Genetic diversity
A 2008 study by Lipinski et al. found that the Birman has one of the lowest levels of genetic diversity of all the breeds studied.
Health
A study in the UK of veterinary records found the Birman to have a life expectancy of 14.39 years based on a sample of 38 cats, higher than the 11.74 average overall.
Paltrinieri, Giraldi, Prolo, Scarpa, et al. (2017) found that Birman cats have a high serum concentration of creatinine and symmetric dimethylarginine, but most Birman cats have higher concentrations of creatinine than SDMA. Creatinine is a creatine phosphate and is produced during metabolism of creatine, and is excreted through urination. SDMA is a methylated form of the amino acid arginine and is released during normal catabolisms of body proteins. Levels of creatinine and SDMA are found when Birman cats are tested for chronic kidney disease, for which they are at high risk. Birman cats are also at risk of developing feline infectious peritonitis; a disease that alters the renal function (creatinine levels in blood and urine) in the cats. In a review of over 5,000 cases of urate urolithiasis the Birman was over-represented with an odds ratio of 6.77.
Feline audiogenic reflex seizures (FARS), a recently discovered type of epilepsy in cats, is believed to be particularly common in Birman cats.
Birman naming conventions
Many Birman breeders follow the French tradition of assigning all kittens born in a particular year given names that begin with the same letter of the alphabet. Countries with breeders using this convention include Australia, Canada, France, New Zealand, the U.K., and the U.S. Kittens born in 2016 would start with 'N', and in 2017 'O', and so on.
Famous Birman cats
Choupette (born 2011), pet of fashion designer Karl Lagerfeld
| Biology and health sciences | Cats | Animals |
146072 | https://en.wikipedia.org/wiki/Stress%20%28biology%29 | Stress (biology) | Stress, whether physiological, biological or psychological, is an organism's response to a stressor such as an environmental condition. When stressed by stimuli that alter an organism's environment, multiple systems respond across the body. In humans and most mammals, the autonomic nervous system and hypothalamic-pituitary-adrenal (HPA) axis are the two major systems that respond to stress. Two well-known hormones that humans produce during stressful situations are adrenaline and cortisol.
The sympathoadrenal medullary (SAM) axis may activate the fight-or-flight response through the sympathetic nervous system, which dedicates energy to more relevant bodily systems to acute adaptation to stress, while the parasympathetic nervous system returns the body to homeostasis.
The second major physiological stress-response center, the HPA axis, regulates the release of cortisol, which influences many bodily functions such as metabolic, psychological and immunological functions. The SAM and HPA axes are regulated by several brain regions, including the limbic system, prefrontal cortex, amygdala, hypothalamus, and stria terminalis. Through these mechanisms, stress can alter memory functions, reward, immune function, metabolism and susceptibility to diseases.
Disease risk is particularly pertinent to mental illnesses, whereby chronic or severe stress remains a common risk factor for several mental illnesses.
Psychology
Acute stressful situations where the stress experienced is severe is a cause of change psychologically to the detriment of the well-being of the individual, such that symptomatic derealization and depersonalization, and anxiety and hyperarousal, are experienced. The International Classification of Diseases includes a group of mental and behavioral disorders which have their aetiology in reaction to severe stress and the consequent adaptive response. Chronic stress, and a lack of coping resources available, or used by an individual, can often lead to the development of psychological issues such as delusions, depression and anxiety (see below for further information). Chronic stress also causes brain atrophy, which is the loss of neurons and the connections between them. It affects the part of the brain that is important for learning, responding to the stressors and cognitive flexibility.
Chronic stressors may not be as intense as acute stressors such as natural disaster or a major accident, but persist over longer periods of time and tend to have a more negative effect on health because they are sustained and thus require the body's physiological response to occur daily. This depletes the body's energy more quickly and usually occurs over long periods of time, especially when these microstressors cannot be avoided (i.e. stress of living in a dangerous neighborhood). See allostatic load for further discussion of the biological process by which chronic stress may affect the body. For example, studies have found that caregivers, particularly those of dementia patients, have higher levels of depression and slightly worse physical health than non-caregivers.
When humans are under chronic stress, permanent changes in their physiological, emotional, and behavioral responses may occur. Chronic stress can include events such as caring for a spouse with dementia, or may result from brief focal events that have long term effects, such as experiencing a sexual assault. Studies have also shown that psychological stress may directly contribute to the disproportionately high rates of coronary heart disease morbidity and mortality and its etiologic risk factors. Specifically, acute and chronic stress have been shown to raise serum lipids and are associated with clinical coronary events.
However, it is possible for individuals to exhibit hardiness—a term referring to the ability to be both chronically stressed and healthy. Even though psychological stress is often connected with illness or disease, most healthy individuals can still remain disease-free after being confronted with chronic stressful events. This suggests that there are individual differences in vulnerability to the potential pathogenic effects of stress; individual differences in vulnerability arise due to both genetic and psychological factors. In addition, the age at which the stress is experienced can dictate its effect on health. Research suggests chronic stress at a young age can have lifelong effects on the biological, psychological, and behavioral responses to stress later in life.
Etymology and historical usage
The term "stress" had none of its contemporary connotations before the 1920s. It is a form of the Middle English destresse, derived via Old French from the Latin stringere, "to draw tight". The word had long been in use in physics to refer to the internal distribution of a force exerted on a material body, resulting in strain. In the 1920s and '30s, biological and psychological circles occasionally used "stress" to refer to a physiological or environmental perturbation that could cause physiological and mental "strain". The amount of strain in reaction to stress depends on the resilience. Excessive strain would appear as illness.
Walter Cannon used it in 1926 to refer to external factors that disrupted what he called homeostasis. But "...stress as an explanation of lived experience is absent from both lay and expert life narratives before the 1930s". Physiological stress represents a wide range of physical responses that occur as a direct effect of a stressor causing an upset in the homeostasis of the body. Upon immediate disruption of either psychological or physical equilibrium the body responds by stimulating the nervous, endocrine, and immune systems. The reaction of these systems causes a number of physical changes that have both short- and long-term effects on the body.
The Holmes and Rahe stress scale was developed as a method of assessing the risk of disease from life changes. The scale lists both positive and negative changes that elicit stress. These include things such as a major holiday or marriage, or death of a spouse and firing from a job.
Biological need for equilibrium
Homeostasis is a concept central to the idea of stress. In biology, most biochemical processes strive to maintain equilibrium (homeostasis), a steady state that exists more as an ideal and less as an achievable condition. Environmental factors, internal or external stimuli, continually disrupt homeostasis; an organism's present condition is a state of constant flux moving about a homeostatic point that is that organism's optimal condition for living. Factors causing an organism's condition to diverge too far from homeostasis can be experienced as stress. A life-threatening situation such as a major physical trauma or prolonged starvation can greatly disrupt homeostasis. On the other hand, an organism's attempt at restoring conditions back to or near homeostasis, often consuming energy and natural resources, can also be interpreted as stress. The brain cannot sustain an equilibrium under chronic stress; the accumulation of such an ever-deepening deficit is called chronic stress.
The ambiguity in defining this phenomenon was first recognized by Hans Selye (1907–1982) in 1926. In 1951 a commentator loosely summarized Selye's view of stress as something that "...in addition to being itself, was also the cause of itself, and the result of itself".
First to use the term in a biological context, Selye continued to define stress as "the non-specific response of the body to any demand placed upon it". Neuroscientists such as Bruce McEwen and Jaap Koolhaas believe that stress, based on years of empirical research, "should be restricted to conditions where an environmental demand exceeds the natural regulatory capacity of an organism". The brain cannot live in an harsh family environment, it needs some sort of stability between another brain. People who have reported being raised in harsh environments such as verbal and physical aggression have showed a more immune dysfunction and more metabolic dysfunction. Indeed, in 1995 Toates already defined stress as a "chronic state that arises only when defense mechanisms are either being chronically stretched or are actually failing," while according to Ursin (1988) stress results from an inconsistency between expected events ("set value") and perceived events ("actual value") that cannot be resolved satisfactorily, which also puts stress into the broader context of cognitive-consistency theory.
Biological background
Stress can have many profound effects on the human biological systems. Biology primarily attempts to explain major concepts of stress using a stimulus-response paradigm, broadly comparable to how a psychobiological sensory system operates. The central nervous system (brain and spinal cord) plays a crucial role in the body's stress-related mechanisms. Whether one should interpret these mechanisms as the body's response to a stressor or embody the act of stress itself is part of the ambiguity in defining what exactly stress is.
The central nervous system works closely with the body's endocrine system to regulate these mechanisms. The sympathetic nervous system becomes primarily active during a stress response, regulating many of the body's physiological functions in ways that ought to make an organism more adaptive to its environment. Below there follows a brief biological background of neuroanatomy and neurochemistry and how they relate to stress.
Stress, either severe, acute stress or chronic low-grade stress may induce abnormalities in three principal regulatory systems in the body: serotonin systems, catecholamine systems, and the hypothalamic-pituitary-adrenocortical axis. Aggressive behavior has also been associated with abnormalities in these systems.
Biology of stress
The brain endocrine interactions are relevant in the translation of stress into physiological and psychological changes. The autonomic nervous system (ANS), as mentioned above, plays an important role in translating stress into a response. The ANS responds reflexively to both physical stressors (for example baroreception), and to higher level inputs from the brain.
The ANS is composed of the parasympathetic nervous system and sympathetic nervous system, two branches that are both tonically active with opposing activities. The ANS directly innervates tissue through the postganglionic nerves, which is controlled by preganglionic neurons originating in the intermediolateral cell column. The ANS receives inputs from the medulla, hypothalamus, limbic system, prefrontal cortex, midbrain and monoamine nuclei.
The activity of the sympathetic nervous system drives what is called the "fight or flight" response. The fight or flight response to emergency or stress involves mydriasis, increased heart rate and force contraction, vasoconstriction, bronchodilation, glycogenolysis, gluconeogenesis, lipolysis, sweating, decreased motility of the digestive system, secretion of the epinephrine and cortisol from the adrenal medulla, and relaxation of the bladder wall. The parasympathetic nervous response, "rest and digest", involves return to maintaining homeostasis, and involves miosis, bronchoconstriction, increased activity of the digestive system, and contraction of the bladder walls. Complex relationships between protective and vulnerability factors on the effect of childhood home stress on psychological illness, cardiovascular illness and adaption have been observed. ANS related mechanisms are thought to contribute to increased risk of cardiovascular disease after major stressful events.
The HPA axis is a neuroendocrine system that mediates a stress response. Neurons in the hypothalamus, particularly the paraventricular nucleus, release vasopressin and corticotropin releasing hormone, which travel through the hypophysial portal vessel where they travel to and bind to the corticotropin-releasing hormone receptor on the anterior pituitary gland. Multiple CRH peptides have been identified, and receptors have been identified on multiple areas of the brain, including the amygdala. CRH is the main regulatory molecule of the release of ACTH.
The secretion of ACTH into systemic circulation allows it to bind to and activate Melanocortin receptor, where it stimulates the release of steroid hormones. Steroid hormones bind to glucocorticoid receptors in the brain, providing negative feedback by reducing ACTH release. Some evidence supports a second long term feedback that is non-sensitive to cortisol secretion. The PVN of the hypothalamus receives inputs from the nucleus of the solitary tract, and lamina terminalis. Through these inputs, it receives and can respond to changes in blood.
The PVN innervation from the brain stem nuclei, particularly the noradrenergic nuclei stimulate CRH release. Other regions of the hypothalamus both directly and indirectly inhibit HPA axis activity. Hypothalamic neurons involved in regulating energy balance also influence HPA axis activity through the release of neurotransmitters such as neuropeptide Y, which stimulates HPA axis activity. Generally, the amygdala stimulates, and the prefrontal cortex and hippocampus attenuate, HPA axis activity; however, complex relationships do exist between the regions.
The immune system may be heavily influenced by stress. The sympathetic nervous system innervates various immunological structures, such as bone marrow and the spleen, allowing for it to regulate immune function. The adrenergic substances released by the sympathetic nervous system can also bind to and influence various immunological cells, further providing a connection between the systems. The HPA axis ultimately results in the release of cortisol, which generally has immunosuppressive effects. However, the effect of stress on the immune system is disputed, and various models have been proposed in an attempt to account for both the supposedly "immunodeficiency" linked diseases and diseases involving hyper activation of the immune system. One model proposed to account for this suggests a push towards an imbalance of cellular immunity(Th1) and humoral immunity(Th2). The proposed imbalance involved hyperactivity of the Th2 system leading to some forms of immune hypersensitivity, while also increasing risk of some illnesses associated with decreased immune system function, such as infection and cancer.
Effects of chronic stress
Chronic stress is a term sometimes used to differentiate it from acute stress. Definitions differ, and may be along the lines of continual activation of the stress response, stress that causes an allostatic shift in bodily functions, or just as "prolonged stress". For example, results of one study demonstrated that individuals who reported relationship conflict lasting one month or longer have a greater risk of developing illness and show slower wound healing. It can also reduce the benefits of receiving common vaccines. Similarly, the effects that acute stressors have on the immune system may be increased when there is perceived stress and/or anxiety due to other events. For example, students who are taking exams show weaker immune responses if they also report stress due to daily hassles. While responses to acute stressors typically do not impose a health burden on young, healthy individuals, chronic stress in older or unhealthy individuals may have long-term effects that are detrimental to health.
Immunological
Acute time-limited stressors, or stressors that lasted less than two hours, results in an up regulation of natural immunity and down regulation of specific immunity. This type of stress saw in increase in granulocytes, natural killer cells, IgA, Interleukin 6, and an increase in cell cytotoxicity. Brief naturalistic stressors elicit a shift from Th1 (cellular) to Th2 (humoral) immunity, while decreased T-cell proliferation, and natural killer cell cytotoxicity. Stressful event sequences did not elicit a consistent immune response; however, some observations such as decreased T-Cell proliferation and cytotoxicity, increase or decrease in natural killer cell cytotoxicity, and an increase in mitogen PHA. Chronic stress elicited a shift toward Th2 immunity, as well as decreased interleukin 2, T cell proliferation, and antibody response to the influenza vaccine. Distant stressors did not consistently elicit a change in immune function. Another response to high impacts of chronic stress that lasts for a long period of time, is more immune dysfunction and more metabolic dysfunction. It is proven in studies that when continuously being in stressful situations, it is more likely to get sick. Also, when being exposed to stress, some claim that the body metabolizes the food in a certain way that adds extra calories to the meal, regardless of the nutritional values of the food.
Infectious
Some studies have observed increased risk of upper respiratory tract infection during chronic life stress. In patients with HIV, increased life stress and cortisol was associated with poorer progression of HIV. Also with an increased level of stress, studies have proven evidence that it can reactivate latent herpes viruses.
Chronic disease
A link has been suggested between chronic stress and cardiovascular disease. Stress appears to play a role in hypertension, and may further predispose people to other conditions associated with hypertension. Stress may precipitate abuse of drugs and/or alcohol. Stress may also contribute to aging and chronic diseases in aging, such as depression and metabolic disorders.
The immune system also plays a role in stress and the early stages of wound healing. It is responsible for preparing the tissue for repair and promoting recruitment of certain cells to the wound area. Consistent with the fact that stress alters the production of cytokines, Graham et al. found that chronic stress associated with care giving for a person with Alzheimer's disease leads to delayed wound healing. Results indicated that biopsy wounds healed 25% more slowly in the chronically stressed group, or those caring for a person with Alzheimer's disease.
Development
Chronic stress has also been shown to impair developmental growth in children by lowering the pituitary gland's production of growth hormone, as in children associated with a home environment involving serious marital discord, alcoholism, or child abuse. Chronic stress also has a lot of illnesses and health care problems other than mental that comes with it. Severe chronic stress for long periods of time can lead to an increased chance of catching illnesses such as diabetes, cancer, depression, heart disease and Alzheimer's disease. More generally, prenatal life, infancy, childhood, and adolescence are critical periods in which the vulnerability to stressors is particularly high. This can lead to psychiatric and physical diseases which have long term impacts on an individual.
Psychopathology
Chronic stress is seen to affect the parts of the brain where memories are processed through and stored. When people feel stressed, stress hormones get over-secreted, which affects the brain. This secretion is made up of glucocorticoids, including cortisol, which are steroid hormones that the adrenal gland releases, although this can increase storage of flashbulb memories it decreases long-term potentiation (LTP). The hippocampus is important in the brain for storing certain kinds of memories and damage to the hippocampus can cause trouble in storing new memories but old memories, memories stored before the damage, are not lost. Also high cortisol levels can be tied to the deterioration of the hippocampus and decline of memory that many older adults start to experience with age. These mechanisms and processes may therefore contribute to age-related disease, or originate risk for earlier-onset disorders. For instance, extreme stress (e.g. trauma) is a requisite factor to produce stress-related disorders such as post-traumatic stress disorder.
Chronic stress also shifts learning, forming a preference for habit based learning, and decreased task flexibility and spatial working memory, probably through alterations of the dopaminergic systems. Stress may also increase reward associated with food, leading to weight gain and further changes in eating habits. Stress may contribute to various disorders, such as fibromyalgia, chronic fatigue syndrome, depression, as well as other mental illnesses and functional somatic syndromes.
Psychological concepts
Eustress
Selye published in year 1975 a model dividing stress into eustress and distress. Where stress enhances function (physical or mental, such as through strength training or challenging work), it may be considered eustress. Persistent stress that is not resolved through coping or adaptation, deemed distress, may lead to anxiety or withdrawal (depression) behavior.
The difference between experiences that result in eustress and those that result in distress is determined by the disparity between an experience (real or imagined) and personal expectations, and resources to cope with the stress. Alarming experiences, either real or imagined, can trigger a stress response.
Coping
Responses to stress include adaptation, psychological coping such as stress management, anxiety, and depression. Over the long term, distress can lead to diminished health and/or increased propensity to illness; to avoid this, stress must be managed.
Stress management encompasses techniques intended to equip a person with effective coping mechanisms for dealing with psychological stress, with stress defined as a person's physiological response to an internal or external stimulus that triggers the fight-or-flight response. Stress management is effective when a person uses strategies to cope with or alter stressful situations.
There are several ways of coping with stress, such as controlling the source of stress or learning to set limits and to say "no" to some of the demands that bosses or family members may make.
A person's capacity to tolerate the source of stress may be increased by thinking about another topic such as a hobby, listening to music, or spending time in a wilderness.
A way to control stress is first dealing with what is causing the stress if it is something the individual has control over. Other methods to control stress and reduce it can be: to not procrastinate and leave tasks for the last minute, do things you like, exercise, do breathing routines, go out with friends, and take a break. Having support from a loved one also helps a lot in reducing stress.
One study showed that the power of having support from a loved one, or just having social support, lowered stress in individual subjects. Painful shocks were applied to married women's ankles. In some trials women were able to hold their husband's hand, in other trials they held a stranger's hand, and then held no one's hand. When the women were holding their husband's hand, the response was reduced in many brain areas. When holding the stranger's hand the response was reduced a little, but not as much as when they were holding their husband's hand. Social support helps reduce stress and even more so if the support is from a loved one.
Cognitive appraisal
Lazarus argued that, in order for a psychosocial situation to be stressful, it must be appraised as such. He argued that cognitive processes of appraisal are central in determining whether a situation is potentially threatening, constitutes a harm/loss or a challenge, or is benign.
Both personal and environmental factors influence this primary appraisal, which then triggers the selection of coping processes. Problem-focused coping is directed at managing the problem, whereas emotion-focused coping processes are directed at managing the negative emotions. Secondary appraisal refers to the evaluation of the resources available to cope with the problem, and may alter the primary appraisal.
In other words, primary appraisal includes the perception of how stressful the problem is and the secondary appraisal of estimating whether one has more than or less than adequate resources to deal with the problem that affects the overall appraisal of stressfulness. Further, coping is flexible in that, in general, the individual examines the effectiveness of the coping on the situation; if it is not having the desired effect, they will, in general, try different strategies.
Assessment
Health risk factors
Both negative and positive stressors can lead to stress. The intensity and duration of stress changes depending on the circumstances and emotional condition of the person with it (Arnold. E and Boggs. K. 2007). Some common categories and examples of stressors include:
Sensory input such as pain, bright light, noise, temperatures, or environmental issues such as a lack of control over environmental circumstances, such as food, air and/or water quality, housing, health, freedom, or mobility.
Social issues can also cause stress, such as struggles with conspecific or difficult individuals and social defeat, or relationship conflict, deception, or break ups, and major events such as birth and deaths, marriage, and divorce.
Life experiences such as poverty, unemployment, clinical depression, obsessive compulsive disorder, heavy drinking, or insufficient sleep can also cause stress. Students and workers may face performance pressure stress from exams and project deadlines.
Adverse experiences during development (e.g. prenatal exposure to maternal stress, poor attachment histories, sexual abuse) are thought to contribute to deficits in the maturity of an individual's stress response systems. One evaluation of the different stresses in people's lives is the Holmes and Rahe stress scale.
General adaptation syndrome
Physiologists define stress as how the body reacts to a stressor - a stimulus, real or imagined. Acute stressors affect an organism in the short term; chronic stressors over the longer term. The general adaptation syndrome (GAS), developed by Hans Selye, is a profile of how organisms respond to stress; GAS is characterized by three phases: a nonspecific alarm mobilization phase, which promotes sympathetic nervous system activity; a resistance phase, during which the organism makes efforts to cope with the threat; and an exhaustion phase, which occurs if the organism fails to overcome the threat and depletes its physiological resources.
Stage 1
Alarm is the first stage, which is divided into two phases: the shock phase and the antishock phase.
Shock phase: During this phase, the body can endure changes such as hypovolemia, hypoosmolarity, hyponatremia, hypochloremia, hypoglycemia—the stressor effect. This phase resembles Addison's disease. The organism's resistance to the stressor drops temporarily below the normal range and some level of shock (e.g. circulatory shock) may be experienced.
Antishock phase: When the threat or stressor is identified or realized, the body starts to respond and is in a state of alarm. During this stage, the locus coeruleus and sympathetic nervous system activate the production of catecholamines including adrenaline, engaging the popularly-known fight-or-flight response. Adrenaline temporarily provides increased muscular tonus, increased blood pressure due to peripheral vasoconstriction and tachycardia, and increased glucose in blood. There is also some activation of the HPA axis, producing glucocorticoids (cortisol, aka the S-hormone or stress-hormone).
Stage 2
Resistance is the second stage. During this stage, increased secretion of glucocorticoids intensifies the body's systemic response. Glucocorticoids can increase the concentration of glucose, fat, and amino acid in blood. In high doses, one glucocorticoid, cortisol, begins to act similarly to a mineralocorticoid (aldosterone) and brings the body to a state similar to hyperaldosteronism. If the stressor persists, it becomes necessary to attempt some means of coping with the stress. The body attempts to respond to stressful stimuli, but after prolonged activation, the body's chemical resources will be gradually depleted, leading to the final stage.
Stage 3
The third stage could be either exhaustion or recovery:
Recovery stage follows when the system's compensation mechanisms have successfully overcome the stressor effect (or have completely eliminated the factor which caused the stress). The high glucose, fat and amino acid levels in blood prove useful for anabolic reactions, restoration of homeostasis and regeneration of cells.
Exhaustion is the alternative third stage in the GAS model. At this point, all of the body's resources are eventually depleted and the body is unable to maintain normal function. The initial autonomic nervous system symptoms may reappear (panic attacks, muscle aches, sore eyes, difficulty breathing, fatigue, heartburn, high blood pressure, and difficulty sleeping, etc.). If stage three is extended, long-term damage may result (prolonged vasoconstriction results in ischemia which in turn leads to cell necrosis), as the body's immune system becomes exhausted, and bodily functions become impaired, resulting in decompensation.
The result can manifest itself in obvious illnesses, such as general trouble with the digestive system (e.g. occult bleeding, melena, constipation/obstipation), diabetes, or even cardiovascular problems (angina pectoris), along with clinical depression and other mental illnesses.
History in research
The current usage of the word stress arose out of Hans Selye's 1930s experiments. He started to use the term to refer not just to the agent but to the state of the organism as it responded and adapted to the environment. His theories of a universal non-specific stress response attracted great interest and contention in academic physiology and he undertook extensive research programs and publication efforts.
While the work attracted continued support from advocates of psychosomatic medicine, many in experimental physiology concluded that his concepts were too vague and unmeasurable. During the 1950s, Selye turned away from the laboratory to promote his concept through popular books and lecture tours. He wrote for both non-academic physicians and, in an international bestseller entitled Stress of Life, for the general public.
A broad biopsychosocial concept of stress and adaptation offered the promise of helping everyone achieve health and happiness by successfully responding to changing global challenges and the problems of modern civilization. Selye coined the term "eustress" for positive stress, by contrast to distress. He argued that all people have a natural urge and need to work for their own benefit, a message that found favor with industrialists and governments. He also coined the term stressor to refer to the causative event or stimulus, as opposed to the resulting state of stress.
Selye was in contact with the tobacco industry from 1958 and they were undeclared allies in litigation and the promotion of the concept of stress, clouding the link between smoking and cancer, and portraying smoking as a "diversion", or in Selye's concept a "deviation", from environmental stress.
From the late 1960s, academic psychologists started to adopt Selye's concept; they sought to quantify "life stress" by scoring "significant life events", and a large amount of research was undertaken to examine links between stress and disease of all kinds. By the late 1970s, stress had become the medical area of greatest concern to the general population, and more basic research was called for to better address the issue. There was also renewed laboratory research into the neuroendocrine, molecular, and immunological bases of stress, conceived as a useful heuristic not necessarily tied to Selye's original hypotheses. The US military became a key center of stress research, attempting to understand and reduce combat neurosis and psychiatric casualties.
The psychiatric diagnosis post-traumatic stress disorder (PTSD) was coined in the mid-1970s, in part through the efforts of anti-Vietnam War activists and the Vietnam Veterans Against the War, and Chaim F. Shatan. The condition was added to the Diagnostic and Statistical Manual of Mental Disorders as posttraumatic stress disorder in 1980. PTSD was considered a severe and ongoing emotional reaction to an extreme psychological trauma, and as such often associated with soldiers, police officers, and other emergency personnel. The stressor may involve threat to life (or viewing the actual death of someone else), serious physical injury, or threat to physical or psychological integrity. In some cases, it can also be from profound psychological and emotional trauma, apart from any actual physical harm or threat. Often, however, the two are combined.
By the 1990s, "stress" had become an integral part of modern scientific understanding in all areas of physiology and human functioning, and one of the great metaphors of Western life. Focus grew on stress in certain settings, such as workplace stress, and stress management techniques were developed. The term also became a euphemism, a way of referring to problems and eliciting sympathy without being explicitly confessional, just "stressed out". It came to cover a huge range of phenomena from mild irritation to the kind of severe problems that might result in a real breakdown of health. In popular usage, almost any event or situation between these extremes could be described as stressful.
The American Psychological Association's 2015 Stress In America Study found that nationwide stress is on the rise and that the three leading sources of stress were "money", "family responsibility", and "work".
| Biology and health sciences | Health and fitness | null |
146103 | https://en.wikipedia.org/wiki/Nonlinear%20system | Nonlinear system | In mathematics and science, a nonlinear system (or a non-linear system) is a system in which the change of the output is not proportional to the change of the input. Nonlinear problems are of interest to engineers, biologists, physicists, mathematicians, and many other scientists since most systems are inherently nonlinear in nature. Nonlinear dynamical systems, describing changes in variables over time, may appear chaotic, unpredictable, or counterintuitive, contrasting with much simpler linear systems.
Typically, the behavior of a nonlinear system is described in mathematics by a nonlinear system of equations, which is a set of simultaneous equations in which the unknowns (or the unknown functions in the case of differential equations) appear as variables of a polynomial of degree higher than one or in the argument of a function which is not a polynomial of degree one.
In other words, in a nonlinear system of equations, the equation(s) to be solved cannot be written as a linear combination of the unknown variables or functions that appear in them. Systems can be defined as nonlinear, regardless of whether known linear functions appear in the equations. In particular, a differential equation is linear if it is linear in terms of the unknown function and its derivatives, even if nonlinear in terms of the other variables appearing in it.
As nonlinear dynamical equations are difficult to solve, nonlinear systems are commonly approximated by linear equations (linearization). This works well up to some accuracy and some range for the input values, but some interesting phenomena such as solitons, chaos, and singularities are hidden by linearization. It follows that some aspects of the dynamic behavior of a nonlinear system can appear to be counterintuitive, unpredictable or even chaotic. Although such chaotic behavior may resemble random behavior, it is in fact not random. For example, some aspects of the weather are seen to be chaotic, where simple changes in one part of the system produce complex effects throughout. This nonlinearity is one of the reasons why accurate long-term forecasts are impossible with current technology.
Some authors use the term nonlinear science for the study of nonlinear systems. This term is disputed by others:
Definition
In mathematics, a linear map (or linear function) is one which satisfies both of the following properties:
Additivity or superposition principle:
Homogeneity:
Additivity implies homogeneity for any rational α, and, for continuous functions, for any real α. For a complex α, homogeneity does not follow from additivity. For example, an antilinear map is additive but not homogeneous. The conditions of additivity and homogeneity are often combined in the superposition principle
An equation written as
is called linear if is a linear map (as defined above) and nonlinear otherwise. The equation is called homogeneous if and is a homogeneous function.
The definition is very general in that can be any sensible mathematical object (number, vector, function, etc.), and the function can literally be any mapping, including integration or differentiation with associated constraints (such as boundary values). If contains differentiation with respect to , the result will be a differential equation.
Nonlinear systems of equations
A nonlinear system of equations consists of a set of equations in several variables such that at least one of them is not a linear equation.
For a single equation of the form many methods have been designed; see Root-finding algorithm. In the case where is a polynomial, one has a polynomial equation such as
The general root-finding algorithms apply to polynomial roots, but, generally they do not find all the roots, and when they fail to find a root, this does not imply that there is no roots. Specific methods for polynomials allow finding all roots or the real roots; see real-root isolation.
Solving systems of polynomial equations, that is finding the common zeros of a set of several polynomials in several variables is a difficult problem for which elaborated algorithms have been designed, such as Gröbner base algorithms.
For the general case of system of equations formed by equating to zero several differentiable functions, the main method is Newton's method and its variants. Generally they may provide a solution, but do not provide any information on the number of solutions.
Nonlinear recurrence relations
A nonlinear recurrence relation defines successive terms of a sequence as a nonlinear function of preceding terms. Examples of nonlinear recurrence relations are the logistic map and the relations that define the various Hofstadter sequences. Nonlinear discrete models that represent a wide class of nonlinear recurrence relationships include the NARMAX (Nonlinear Autoregressive Moving Average with eXogenous inputs) model and the related nonlinear system identification and analysis procedures. These approaches can be used to study a wide class of complex nonlinear behaviors in the time, frequency, and spatio-temporal domains.
Nonlinear differential equations
A system of differential equations is said to be nonlinear if it is not a system of linear equations. Problems involving nonlinear differential equations are extremely diverse, and methods of solution or analysis are problem dependent. Examples of nonlinear differential equations are the Navier–Stokes equations in fluid dynamics and the Lotka–Volterra equations in biology.
One of the greatest difficulties of nonlinear problems is that it is not generally possible to combine known solutions into new solutions. In linear problems, for example, a family of linearly independent solutions can be used to construct general solutions through the superposition principle. A good example of this is one-dimensional heat transport with Dirichlet boundary conditions, the solution of which can be written as a time-dependent linear combination of sinusoids of differing frequencies; this makes solutions very flexible. It is often possible to find several very specific solutions to nonlinear equations, however the lack of a superposition principle prevents the construction of new solutions.
Ordinary differential equations
First order ordinary differential equations are often exactly solvable by separation of variables, especially for autonomous equations. For example, the nonlinear equation
has as a general solution (and also the special solution corresponding to the limit of the general solution when C tends to infinity). The equation is nonlinear because it may be written as
and the left-hand side of the equation is not a linear function of and its derivatives. Note that if the term were replaced with , the problem would be linear (the exponential decay problem).
Second and higher order ordinary differential equations (more generally, systems of nonlinear equations) rarely yield closed-form solutions, though implicit solutions and solutions involving nonelementary integrals are encountered.
Common methods for the qualitative analysis of nonlinear ordinary differential equations include:
Examination of any conserved quantities, especially in Hamiltonian systems
Examination of dissipative quantities (see Lyapunov function) analogous to conserved quantities
Linearization via Taylor expansion
Change of variables into something easier to study
Bifurcation theory
Perturbation methods (can be applied to algebraic equations too)
Existence of solutions of Finite-Duration, which can happen under specific conditions for some non-linear ordinary differential equations.
Partial differential equations
The most common basic approach to studying nonlinear partial differential equations is to change the variables (or otherwise transform the problem) so that the resulting problem is simpler (possibly linear). Sometimes, the equation may be transformed into one or more ordinary differential equations, as seen in separation of variables, which is always useful whether or not the resulting ordinary differential equation(s) is solvable.
Another common (though less mathematical) tactic, often exploited in fluid and heat mechanics, is to use scale analysis to simplify a general, natural equation in a certain specific boundary value problem. For example, the (very) nonlinear Navier-Stokes equations can be simplified into one linear partial differential equation in the case of transient, laminar, one dimensional flow in a circular pipe; the scale analysis provides conditions under which the flow is laminar and one dimensional and also yields the simplified equation.
Other methods include examining the characteristics and using the methods outlined above for ordinary differential equations.
Pendula
A classic, extensively studied nonlinear problem is the dynamics of a frictionless pendulum under the influence of gravity. Using Lagrangian mechanics, it may be shown that the motion of a pendulum can be described by the dimensionless nonlinear equation
where gravity points "downwards" and is the angle the pendulum forms with its rest position, as shown in the figure at right. One approach to "solving" this equation is to use as an integrating factor, which would eventually yield
which is an implicit solution involving an elliptic integral. This "solution" generally does not have many uses because most of the nature of the solution is hidden in the nonelementary integral (nonelementary unless ).
Another way to approach the problem is to linearize any nonlinearity (the sine function term in this case) at the various points of interest through Taylor expansions. For example, the linearization at , called the small angle approximation, is
since for . This is a simple harmonic oscillator corresponding to oscillations of the pendulum near the bottom of its path. Another linearization would be at , corresponding to the pendulum being straight up:
since for . The solution to this problem involves hyperbolic sinusoids, and note that unlike the small angle approximation, this approximation is unstable, meaning that will usually grow without limit, though bounded solutions are possible. This corresponds to the difficulty of balancing a pendulum upright, it is literally an unstable state.
One more interesting linearization is possible around , around which :
This corresponds to a free fall problem. A very useful qualitative picture of the pendulum's dynamics may be obtained by piecing together such linearizations, as seen in the figure at right. Other techniques may be used to find (exact) phase portraits and approximate periods.
Types of nonlinear dynamic behaviors
Amplitude death – any oscillations present in the system cease due to some kind of interaction with other system or feedback by the same system
Chaos – values of a system cannot be predicted indefinitely far into the future, and fluctuations are aperiodic
Multistability – the presence of two or more stable states
Solitons – self-reinforcing solitary waves
Limit cycles – asymptotic periodic orbits to which destabilized fixed points are attracted.
Self-oscillations – feedback oscillations taking place in open dissipative physical systems.
Examples of nonlinear equations
Algebraic Riccati equation
Ball and beam system
Bellman equation for optimal policy
Boltzmann equation
Colebrook equation
General relativity
Ginzburg–Landau theory
Ishimori equation
Kadomtsev–Petviashvili equation
Korteweg–de Vries equation
Landau–Lifshitz–Gilbert equation
Liénard equation
Navier–Stokes equations of fluid dynamics
Nonlinear optics
Nonlinear Schrödinger equation
Power-flow study
Richards equation for unsaturated water flow
Self-balancing unicycle
Sine-Gordon equation
Van der Pol oscillator
Vlasov equation
| Mathematics | Dynamical systems | null |
146151 | https://en.wikipedia.org/wiki/Washington%20Metro | Washington Metro | The Washington Metro, often abbreviated as the Metro and formally the Metrorail, is a rapid transit system serving the Washington metropolitan area of the United States. It is administered by the Washington Metropolitan Area Transit Authority (WMATA), which also operates the Metrobus service under the Metro name. Opened in 1976, the network now includes six lines, 98 stations, and of route.
Metro serves Washington, D.C. and the states of Maryland and Virginia. In Maryland, Metro provides service to Montgomery and Prince George's counties; in Virginia, to Arlington, Fairfax and Loudoun counties, and to the independent city of Alexandria. The system's most recent expansion, which is the construction of a new station (and altering the line), serving Potomac Yard, opened on May 19, 2023. It operates mostly as a deep-level subway in more densely populated parts of the D.C. metropolitan area (including most of the District itself), while most of the suburban tracks are at surface level or elevated. The longest single-tier escalator in the Western Hemisphere, spanning , is located at Metro's deep-level station.
In , the system had a ridership of , or about per weekday as of , making it the second-busiest heavy rail rapid transit system in the United States, in number of passenger trips, after the New York City Subway, and the sixth-busiest in North America. In June 2008, Metro set a monthly ridership record with 19,729,641 trips, or 798,456 per weekday. Fares vary based on the distance traveled, the time of day, and the type of card used by the passenger. Riders enter and exit the system using a proximity card called SmarTrip.
History
During the 1950s, plans were laid for a massive freeway system in Washington, D.C. Harland Bartholomew, who chaired the National Capital Planning Commission, thought that a rail transit system would never be self-sufficient because of low-density land uses and general transit ridership decline. But the plan met fierce opposition, and was altered to include a Capital Beltway system plus rail line radials. The Beltway received full funding along with additional funding from the Inner Loop Freeway system project that was partially reallocated toward construction of the Metro system.
In 1960, the federal government created the National Capital Transportation Agency to develop a rapid rail system. In 1966, a bill creating WMATA was passed by the federal government, the District of Columbia, Virginia, and Maryland, with planning power for the system being transferred to it from the NCTA. An early proposal map from 1967 was more extensive than what was ultimately approved, with the Red Line's western terminus being in Germantown instead of Shady Grove.
WMATA approved plans for a regional system on March 1, 1968. The plan consisted of a core regional system, which included the original five Metro lines, as well as several future extensions, many of which were not constructed. The first experimental Metro station was built above ground in May 1968 for a cost of $69,000. It was and meant to test construction techniques, lighting, and acoustics before full-scale construction efforts.
Construction began after a groundbreaking ceremony on December 9, 1969, when Secretary of Transportation John A. Volpe, District Mayor Walter Washington, and Maryland Governor Marvin Mandel tossed the first spade of dirt at Judiciary Square.
The first portion of the system opened March 27, 1976, with available on the Red Line with five stations from to , all in Washington, D.C. All rides were free that day, with the first train departing the Rhode Island Avenue stop with Metro officials and special guests, and the second with members of the general public. Arlington County, Virginia was linked to the system on July 1, 1977; Montgomery County, Maryland, on February 6, 1978; Prince George's County, Maryland, on November 17, 1978; and Fairfax County, Virginia, and Alexandria, Virginia, on December 17, 1983. Metro reached Loudoun County on November 15, 2022. Underground stations were built with cathedral-like arches of concrete, highlighted by soft, indirect lighting. The name Metro was suggested by Massimo Vignelli, who designed the signage for the system as well as for the New York City Subway.
The , 83-station system was completed with the opening of the Green Line segment to on January 13, 2001. However, this did not mean the end of the system's growth. A extension of the Blue Line to and opened on December 18, 2004. The first infill station, New York Ave–Florida Ave–Gallaudet University (now ) on the Red Line between and , opened on November 20, 2004. Construction began in March 2009 for an extension to Dulles Airport to be built in two phases. The first phase, five stations connecting East Falls Church to Tysons Corner and Wiehle Avenue in Reston, opened on July 26, 2014. The second phase to Ashburn opened November 15, 2022, after many delays. The second infill station, on the Blue and Yellow Lines between and , opened on May 19, 2023.
Metro construction required billions of federal dollars, originally provided by Congress under the authority of the National Capital Transportation Act of 1969. The cost was paid with 67% federal money and 33% local money. This act was amended on January 3, 1980, by the National Capital Transportation Amendment of 1979 (also known as the Stark-Harris Act), which authorized additional funding of $1.7 billion to permit the completion of of the system as provided under the terms of a full funding grant agreement executed with WMATA in July 1986, which required 20% to be paid from local funds. On November 15, 1990, the National Capital Transportation Amendments of 1990 authorized an additional $1.3 billion in federal funds for construction of the remaining of the system, completed via the execution of full funding grant agreements, with a 63% federal/37% local matching ratio.
In February 2006, Metro officials chose Randi Miller, a car dealership employee from Woodbridge, Virginia, to record new "doors opening", "doors closing", and "please stand clear of the doors, thank you" announcements after winning an open contest to replace the messages recorded by Sandy Carroll in 1996. The "Doors Closing" contest attracted 1,259 contestants from across the country.
Over the years, a lack of investment in Metro caused it to break down, and there have been several fatal incidents on the Washington Metro due to mismanagement and broken-down infrastructure. By 2016, according to The Washington Post, on-time rates had dropped to 84%, and Metro service was frequently disrupted during rush hours because of a combination of equipment, rolling stock, track, and signal malfunctions. WMATA did not receive dedicated funding from the three jurisdictions it served, Maryland, Virginia, and D.C., until 2018.
Seeking to address negative perceptions of its performance, in 2016, WMATA announced an initiative called "Back2Good," focusing on addressing a wide array of rider concerns, from improving safety to adding Internet access to stations and train tunnels.
In May 2018, Metro announced an extensive renovation of platforms at 20 stations across the system, spanning all lines except the Silver Line. The Blue and Yellow Lines south of were closed from May 25 to September 9, 2019, in what would be the longest line closure in Metro's history. Additional stations would be repaired between 2020 and 2022, but the corresponding lines would not be closed completely. The project would cost $300 to $400 million and would be Metro's first major project since its construction.
In March 2022, Metro announced that beginning on September 10, 2022, it would suspend all service on the Yellow Line for seven to eight months to complete repairs and rebuilding work on its bridge over the Potomac River and its tunnel leading into the station at . Metro stated that this was the first significant work that the tunnel and bridge had undergone since they were first constructed over forty years prior. Service on the Yellow Line resumed on May 7, 2023, but with its northeastern terminus truncated from to .
Opening dates
The following is a list of opening dates for track segments and infill stations on the Washington Metro. The entries in the "from" and "to" columns correspond to the boundaries of the extension or station that opened on the specified date, not to the lines' terminals.
Rush+ and late-night service patterns
On December 31, 2006, an 18-month pilot program began to extend service on the Yellow Line to Fort Totten over existing Green Line trackage. This extension was later made permanent. Starting June 18, 2012, the Yellow Line was extended again along existing track as part of the Rush+ program, with an extension to Greenbelt on the northern end and with several trains diverted to Franconia–Springfield on the southern end. These Rush+ extensions were discontinued on June 25, 2017.
In addition to expanding the system, Metro expanded the operating hours over the first 40 years. Though it originally opened with weekday-only service from 6 a.m. to 8 p.m, financial paperwork assumed prior to opening that it would eventually operate from 5 a.m. to 1 a.m. seven days a week. It never operated exactly on that schedule but the hours did expand, sometimes beyond that. On September 25, 1978, Metro extended its weekday closing time from 8 p.m. to midnight and 5 days later it started Saturday service from 8 a.m. to Midnight. Metrorail kicked off Sunday service from 10 a.m. to 6 p.m. on September 2, 1979, and on June 29, 1986, the Sunday closing time was pushed back to midnight. Metro started opening at 5:30 a.m., a half an hour earlier, on weekdays starting on July 1, 1988. On November 5, 1999, weekend service was extended to 1:00 a.m., and on June 30, 2000, it was expanded to 2:00 a.m. On July 5, 2003, weekend hours were extended again with the system opening an hour earlier, at 7:00 a.m. and closing an hour later at 3:00 a.m. On September 27, 2004, Metro again pushed weekday opening time half an hour earlier, this time to 5 a.m.
In 2016, Metro began temporarily scaling back service hours to allow for more maintenance. On June 3, 2016, they ended late-night weekend service with Metrorail closing at midnight. Hours were adjusted again the following year starting on June 25, 2017, with weeknight service ending a half-hour earlier at 11:30 p.m.; Sunday service trimmed to start an hour later – at 8 a.m. – and end an hour early at 11 p.m.; and late-night service partially restored to 1 a.m. The service schedule was approved until June 2019.
On January 29, 2020, Metro announced that it would be activating its pandemic response plans in preparation for the looming COVID-19 pandemic, which would be declared a pandemic by the World Health Organization on March 11. At that time, Metro announced that it would reduce its service hours from 5:00 a.m. to 9:00 p.m. on weekdays and 8:00 a.m. to 9:00 p.m. on weekends beginning on March 16 to accommodate for train cleaning and additional track work. As of 2022, pre-COVID service hours have been restored with pre-2016 Sunday service hours.
Busiest days
The highest ridership for a single day was on the day of the first inauguration of Barack Obama, January 20, 2009, with 1.12 million riders. It broke the previous record, set the day before, of 866,681 riders. June 2008 set several ridership records: the single-month ridership record of 19,729,641 total riders, the record for highest average weekday ridership with 1,044,400 weekday trips, had five of the ten highest ridership days, and had 12 weekdays in which ridership exceed 800,000 trips. The Sunday record of 616,324 trips was set on January 18, 2009, during Obama's pre-inaugural events, the day the Obamas arrived in Washington and hosted a concert on the steps of the Lincoln Memorial. It broke the record set on the 4th of July, 1999.
On January 21, 2017, the 2017 Women's March, set an all-time record in Saturday ridership with 1,001,616 trips.
The previous record was set on October 30, 2010, with 825,437 trips during the Rally to Restore Sanity and/or Fear. Prior to 2010, the record had been set on June 8, 1991, at 786,358 trips during the Desert Storm rally.
Architecture
Many Metro stations were designed by Chicago architect Harry Weese and are examples of late 20th century modern architecture. With their heavy use of exposed concrete and repetitive design motifs, Metro stations display aspects of Brutalist design. The stations also reflect the influence of Washington's neoclassical architecture in their overarching coffered ceiling vaults. Weese worked with Cambridge, Massachusetts-based lighting designer Bill Lam on the indirect lighting used throughout the system. All of Metro's original Brutalist stations are found in Downtown Washington, D.C., and neighboring urban corridors of Arlington, Virginia, while newer stations incorporate simplified cost-efficient designs.
In 2007, the design of the Metro's vaulted-ceiling stations was voted number 106 on the "America's Favorite Architecture" list compiled by the American Institute of Architects (AIA), and was the only Brutalist design to win a place among the 150 selected by this public survey.
In January 2014, the AIA announced that it would present its Twenty-five Year Award to the Washington Metro system for "an architectural design of enduring significance" that "has stood the test of time by embodying architectural excellence for 25 to 35 years". The announcement cited the key role of Weese, who conceived and implemented a "common design kit-of-parts", which continues to guide the construction of new Metro stations over a quarter-century later, albeit with designs modified slightly for cost reasons.
Beginning in 2003, canopies were added to existing exits of underground stations due to the wear and tear seen on escalators due to exposure to the elements.
System
Since opening in 1976, the Metro network has grown to include six lines, 98 stations, and of route. The rail network is designed according to a spoke–hub distribution paradigm, with rail lines running between downtown Washington and its nearby suburbs. The system extensively uses interlining: running more than one service on the same track. There are six operating lines. The system's official map was designed by noted graphic designer Lance Wyman and Bill Cannan while they were partners in the design firm of Wyman & Cannan in New York City.
About of Metro's track is underground, as are 47 of the 98 stations. Track runs underground mostly within the District and high-density suburbs. Surface track accounts for about of the total, and aerial track makes up . The system operates on a track gauge of , which is narrower than but within the tolerance of standard-gauge railways.
Previously, the least time to travel through 97 stations using only mass transit was 8 hours 54 minutes, a record set by travel blogger Lucas Wall on November 16, 2022, the first full day that Phase 2 of the Silver Line was in passenger operation. This record was broken by a student named Claire Aguayo, who did it in 8 hours and 36 minutes on January 23, 2023. Both of these runs were before the station opened on May 19, 2023, making them no longer current.
To gain revenues, WMATA has started to allow retail ventures in Metro stations. WMATA has authorized DVD-rental vending machines and ticket booths for the Old Town Trolley Tours and is seeking additional retail tenants.
Financing
Metro relies extensively on passenger fares and appropriated financing from the Maryland, Virginia, and Washington D.C., governments, which are represented on Metro's board of directors. In 2018, Maryland, Virginia and Washington, D.C., agreed to contribute $500 million annually to Metro's capital budget. Until then, the system did not have a dedicated revenue stream as other cities' mass transit systems do. Critics allege that this has contributed to Metro's recent history of maintenance and safety problems.
For Fiscal Year 2019, the estimated farebox recovery ratio (fare revenue divided by operating expenses) was 62 percent, based on the WMATA-approved budget.
Infrastructure
Stations
There are 40 stations in the District of Columbia, 15 in Prince George's County, 13 in Fairfax County, 11 in Montgomery County, 11 in Arlington County, 5 in the City of Alexandria, and 3 in Loudoun County. The most recent station was opened on May 19, 2023, an infill station at . At below the surface, the station on the Red Line is the deepest in the system. There are no escalators; high-speed elevators take 20 seconds to travel from the street to the station platform. The station, one stop to the north of the Forest Glen station, has the longest continuous escalator in the US and in the Western Hemisphere, at . The station is the deepest station on the Orange/Blue/Silver Line, at below street level. The station features the second-longest continuous escalator in the Metro system at ; an escalator ride between the street and mezzanine levels takes nearly two minutes.
The system is not centered on any single station, but is at the intersection of the Red, Orange, Blue, and Silver Lines. The station is also the location of WMATA's main sales office, which closed in 2022. Metro has designated five other "core stations" that have high passenger volume, including: , transfer station for the Red, Green, and Yellow Lines; , transfer station for the Orange, Blue, Silver, Green, and Yellow Lines; , the busiest station by passenger boardings; ; and .
To deal with the high number of passengers in transfer stations, Metro is studying the possibility of building pedestrian connections between nearby core transfer stations. For example, a passage between Metro Center and Gallery Place stations would allow passengers to transfer between the Orange/Blue/Silver and Yellow/Green Lines without going to one stop on the Red Line or taking a slight detour via L’Enfant Plaza. Another tunnel between Farragut West and Farragut North stations would allow transfers between the Red and Orange/Blue/Silver lines, decreasing transfer demand at Metro Center by an estimated 11%. The Farragut pedestrian tunnel has yet to be physically implemented, but was added in virtual form effective October 28, 2011: the SmarTrip system now interprets an exit from one Farragut station and entrance to the other as part of a single trip, allowing cardholders to transfer on foot without having to pay a second full fare.
Rolling stock
Metro's fleet consists of 1,216 rail cars, each long, with 1,208 in active revenue service as of May 2024. Though operating rules currently limit trains to (except on the Green line, where they can go up to ), all trains have a maximum speed of , and average , including stops. All cars operate as married pairs (consecutively numbered even-odd with a cab at each end of the pair except 7000-series railcars), with systems shared across the pair.
In the "Active railcars" table, font in bold represents the railcars that are currently in service, while the regular font represents cars that are temporarily out of service
Metro's rolling stock was acquired in seven phases, and each version of car is identified with a separate series number.
The original order of 300 railcars (all of which have been retired as of July 1, 2017) was manufactured by Rohr Industries, with final delivery in 1978. These cars are numbered 1000–1299 and were rehabilitated in the mid-1990s.
Breda Costruzioni Ferroviarie (Breda), manufactured the second order of 76 cars delivered in 1983 and 1984. These cars, numbered 2000–2075, were rehabilitated in the early 2000s by Alstom in Hornell, New York. All 2000-series cars were retired by May 10, 2024.
A third order of 290 cars, also from Breda, were delivered between 1984 and 1988. These cars are numbered 3000–3289 and were rehabilitated by Alstom in the mid-2000s.
A fourth order of 100 cars from Breda, numbered 4000–4099, were delivered between 1991 and 1994. All 4000-series cars were retired by July 1, 2017.
A fifth order of 192 cars was manufactured by Construcciones y Auxiliar de Ferrocarriles (CAF) of Spain. These cars are numbered 5000–5191 and were delivered from 2001 through 2004. Most 5000-series cars were retired in October 2018 and the last few in spring 2019.
A sixth order of 184 cars from Alstom Transportation, are numbered 6000–6183 and were delivered between 2005 and 2007. The cars have body shells built in Barcelona, Spain with assembly completed in Hornell, New York.
The 7000-series railcars, built by Kawasaki Heavy Industries Rolling Stock Company of Kobe, Japan, were delivered for on-site testing during winter 2013–2014, and first entered service on April 14, 2015, on the Blue Line. The cars are different from previous models in that while still operating as married pairs, the cab in one car is eliminated, turning it into a B car. This design allows for increased passenger capacity, elimination of redundant equipment, greater energy efficiency, and lower maintenance costs. The National Transportation Safety Board investigation of the fatal June 22, 2009, accident led it to conclude that the 1000-series cars are unsafe and unable to protect passengers in a crash. As a result, on July 26, 2010, Metro voted to purchase 300 7000-series cars, which replaced the remaining 1000-series cars. An additional 128 7000-series cars were also ordered to serve the Silver Line to Dulles Airport (64 for each phase). In April 2013, Metro placed another order for 100 7000-series cars, which replaced all of the 4000-series cars. On July 13, 2015, WMATA used their final option and purchased an additional 220 7000-series railcars for fleet expansion and to replace the 5000-series railcars, bringing the total order number to 748 railcars. On February 26, 2020, WMATA accepted the delivery of the final 7000-series car.
The 8000-series cars will be constructed by Hitachi Rail. While these railcars would have a similar appearance to the 7000-series, the 8000-series would include more features such as "smart doors" that detect obstruction, high-definition security cameras, more space between seats, wider aisles, and non-slip flooring. In September 2018, Metro issued a request for proposals from manufacturers for 256 railcars with options for a total of up to 800. The first order would replace the 2000 and 3000-series equipment, while the options, if selected, would allow the agency to increase capacity and retire the 6000-series.
Signaling and operation
During normal passenger operation on revenue tracks, trains are designed to be controlled by an integrated Automatic Train Operation (ATO) and Automatic Train Control (ATC) system that accelerates and brakes trains automatically without operator intervention. All trains are still staffed with train operators who open and close the doors, make station announcements, and supervise their trains. The system was designed so that an operator could manually operate a train when necessary.
Since June 2009, when two Red Line trains collided and killed nine people due in part to malfunctions in the ATC system, all Metro trains have been manually operated. The current state of manual operation has led to heavily degraded service, with new manual requirements such as absolute blocks, speed restrictions, and end-of-platform stopping leading to increased headways between trains, increased dwell time, and worse on-time performance. Metro originally planned to have all trains be automated again by 2017, but those plans were shelved in early 2017 in order to focus on more pressing safety and infrastructure issues. In March 2023, Metro announced plans to re-automate the system by December of that year, but announced in September that these plans would be delayed until 2024. ATO resumed only on the Red Line on December 15, 2024, and the highest speed is now 75mph, with ATO scheduled to resume on the rest of the system in 2025.
The train doors were originally designed to be opened and closed automatically and the doors would re-open if an object blocked them, much as elevator doors do. Almost immediately after the system opened in 1976 Metro realized these features were not conducive to safe or efficient operation and they were disabled. Metro began testing reinstating automatic train door opening in March 2019, citing delays and potential human error. If a door tries to close and it meets an obstruction, the operator must re-open the door. In October 2023, automatic train door opening, where train doors will automatically open upon alighting, was restored to a limited number of trains on the Red Line. Operators must manually close the doors after they open. WMATA claims that automatic door opening provides a safety benefit since it eliminates potential human error resulting in the doors opening on the wrong side and a reduction in the wait time before doors opening, improving the customer experience and station dwell times.
Hours and headways
Metrorail begins service at 5 am Monday through Friday, 7 am on Saturdays and Sundays; it ends service at midnight Monday through Thursday, 1:00 am Friday and Saturday, and midnight on Sundays, although the last trains leave the end stations inbound about half an hour before these times. Pre-pandemic, trains ran more frequently during rush hours on all lines, with scheduled peak hour headways of 4 minutes on the Red Line and 8 minutes on all other lines. Headways were much longer during midday and evening on weekdays and all day weekends. The midday six-minute headways were based on a combination of two Metrorail lines (Orange/Blue and Yellow/Green) as each route could run every 12 minutes (4 minutes for the Blue/Orange/Silver segment); in the case of the Red Line, every other train bound for Glenmont terminated at Silver Spring instead. Night and weekend service varied between 6 and 20 minutes, with trains generally scheduled only every 15 to 20 minutes.
Other service truncations also occur in the system during rush hour service only. On the Red Line, every other train bound for Shady Grove terminated at until December 2018, in addition to the alternating terminations at Silver Spring mentioned above. For the Yellow Line, all non Rush+ trains bound for and all normal trains bound for terminate at . These are primarily instituted due to a limited supply of rail cars and the locations of pocket tracks throughout the system. However, as of July 2019, both Red Line service truncations have ended, and
as of April 2019, the Yellow Line served Greenbelt at all times. When the Yellow Line reopened on May 7, 2023, following major maintenance work, the Mount Vernon Square turnback was reinstated at all times, which has not happened since 2006.
Until 1999, Metro ended service at midnight every night, and weekend service began at 8 am. That year, WMATA began late-night service on Fridays and Saturdays until 1 am. By 2007, with encouragement from businesses, that closing time had been pushed back to 3 am, with peak fares in effect for entries after midnight. There were plans floated to end late-night service due to costs in 2011, but they were met with resistance by riders.
WMATA temporarily discontinued late night rail service on May 30, 2016, so that Metro can conduct an extensive track rehabilitation program in an effort to improve the system's reliability.
On June 25, 2017, Metro cut its hours of operation with closing at 11:30 PM Monday–Thursday, 1 AM on Friday and Saturday, and 11 PM on Sunday, with the last trains leaving the end stations inbound about half an hour before these times. As of 2022, the pre-2017 service hours have been restored.
Special service patterns
Metro runs special service patterns on holidays and when events in Washington may require additional service. Independence Day activities require Metro to adjust service to provide extra capacity to and from the National Mall. WMATA makes similar adjustments during other events, such as presidential inaugurations. Due to security concerns related to the January 6 United States Capitol attack, several Metro stations were closed for the 2021 Inauguration. Metro has altered service and used some stations as entrances or exits only to help manage congestion.
Rush Plus
In 2012, WMATA announced enhanced rush period service that was implemented on June 18, 2012, under the name "Rush+" (or "Rush Plus"). Rush Plus service occurred only during portions of peak service: 6:30–9:00 AM and 3:30–6:00 PM, Monday through Friday.
The Rush+ realignment was intended to free up space in the Rosslyn Portal (the tunnel between Rosslyn and Foggy Bottom), which operates at full capacity already. When Silver Line service began, those trains would be routed through the tunnel, and so some of what were Blue Line trains to were now diverted across the Fenwick Bridge to become Yellow Line trains running all the way along the Green Line to . Select Yellow Line trains running south diverted along the Blue Line to (as opposed to the normal Yellow line terminus at ). Until the start of Silver Line service, excess Rosslyn Tunnel capacity was used by additional Orange Line trains that traveled along the Blue Line to Largo (as opposed to the normal Orange Line terminus at ). Rush+ had the additional effect of giving some further number of passengers transfer-free journeys, though severely increasing headways for the portion of the Blue Line running between and . In May 2017, Metro announced that Yellow Rush+ service would be eliminated effective June 25, 2017.
COVID-19 and 7000-series derailment (2020–present)
Headways have been lengthened as a result of the COVID-19 pandemic in Washington, D.C., starting early 2020. Near-pre-pandemic service was restored at times until October 2021, but due to the 7000-series derailment near Arlington Cemetery, and subsequent removal of all 7000-series cars from service (which made up 60% of the WMATA fleet), headways were lengthened again to every 15 minutes on the Red Line and every 30 minutes on all other lines beginning October 19, 2021.
Since then, with more 7000-series cars returning, headways have been gradually restored to near-pre-pandemic levels, especially outside of peak times, with ridership also increasing as a result. As of September 2024, several lines are actually more frequent than 2019 levels during certain times of day on weekdays and/or weekends. The Red Line's evening headways improved from every 15 minutes in 2019 to every 10 minutes in 2024. In 2019, all lines except the Red Line had 20-minute evening headways, whereas in 2024 the Green and Yellow Lines run every 8 minutes during evenings and the Blue, Orange, and Silver Lines every 15. Sunday service improved to match Monday-Friday off-peak and Saturday levels of every 6 minutes on the Red Line, every 8 minutes on the Green and Yellow Lines, and every 12 minutes on the Blue, Orange, and Silver Lines, compared to the previous 8 minutes on the Red Line and 15 minutes on all other lines. The Yellow and Green Lines also currently run every 6 minutes during rush hours starting 2023 for the first time since major peak service cuts in 2017 that eliminated Rush Plus and decreased rush hour frequencies on all lines except the Blue Line from 6 to 8 minutes.
Current headways by line
Headways as of December 15, 2024.
Current average headways by line segment
Headways as of November 9, 2024. Calculated using trains per hour and rounded to nearest minute.
Passenger information systems
A passenger information display system (PIDS) was installed in all Metrorail stations in 2000. Displays are located on all track platforms and at the mezzanine entrances of stations. They provide real-time information on next train arrivals, including the line, destination, number of cars in the train, and estimated wait time. The displays also show information about delayed trains, emergency announcements, and other bulletins. The signs were upgraded in 2013 to better reflect Rush Plus and Silver Line schedules, and to prioritize next-train arrival information over other announcements. New digital PIDS signs were installed at the six stations south of National Airport in summer 2019 as part of the Platform Improvement Project.
WMATA also provides current train and related information to customers with conventional web browsers, as well as users of smartphones and other mobile devices. In 2010 Metro began sharing its PIDS data with outside software developers, for use in creating additional real-time applications for mobile devices. Free apps are available to the public on major mobile device software platforms (iOS, Android, Windows Phone, Palm). WMATA also began providing real-time train information by phone in 2010.
Fare structure
Riders enter and exit the system using a stored-value card in the form of a proximity card known as SmarTrip. The fare is deducted from the balance of the card when exiting. SmarTrip cards can be purchased at station vending machines, online or at retail outlets, and can store up to $300 in value. Metro also accepts Baltimore's CharmCard, a similar contactless payment card system.
Metro fares vary based on the distance traveled and the time of day at entry. Fares (effective 2024) range from $2.25 to $6.75, depending on the distance traveled during weekdays prior to 9:30 PM and $2.25 to $2.50 on weekends or after 9:30 PM on weekdays at the time of tapping in. Discounted fares from 50% to 100% are available for DC school children, SNAP Recipients in Maryland, Virginia, and Washington DC, disabled people, and senior citizens.
Parking fees at Metro stations range from $3.00 to $5.20 on weekdays for riders; non-rider fees range from $3.00 to $10.00. Parking is free on Saturdays, Sundays, and federal holidays.
Since June 25, 2017, the first fare hike in three years, peak-period rail fares increased 10 cents, with $2.25 as the new minimum and $6.00 as the maximum one-way fare. Off-peak fares rose 25 cents, to a $2.00 minimum and $3.85 maximum, as will bus fares. A new one-day unlimited rail/bus pass became available for $14.75, which is presently available for $13.50.
On June 24, 2024, WMATA announced another fare hike effective June 30, 2024, with a general increase of 12.5% to most services. Of the fare increases, the rail fare during the weekday increased to range from $2.25 to $6.75, while the flat $2.00 rate during late night (after 9:30) and weekend hours was replaced to range from $2.25 to $2.50 depending on the distance traveled.
Passengers may purchase passes at farecard vending machines. Passes are loaded onto the same SmarTrip cards as stored value, but grant riders unlimited travel within the system for a certain period of time. The period of validity starts with the first use. Four types of passes are currently sold:
A 1-Day Unlimited Pass for $13.50, valid for one day of unlimited Metrorail and Metrobus travel. The pass expires at the end of the operating day.
A 3-Day Unlimited Pass for $33.75, valid for three consecutive days of unlimited Metrorail and Metrobus travel.
A 7-Day Short Trip Unlimited Pass for $40.50, valid for seven consecutive days for Metrorail trips costing up to $4.50. If the trip costs more than $4.50, the difference is deducted from the cash balance of a SmarTrip card, possibly after the necessary value is added at the Exitfare machine. A non-negative stored value is required to enter and exit the Metrorail system.
A 7-Day Unlimited Pass for $60.75, valid for seven consecutive days of unlimited Metrorail and Metrobus travel.
In addition, Metro sells the Monthly Unlimited Pass, formerly called SelectPass, available for purchase online only by registered SmarTrip cardholders, valid for trips up to a specified value for a specific calendar month, with the balance being deducted from the card's cash value similarly to the Short Trip Pass. The pass is priced based on 18 days of round-trip travel.
Users can add value to any farecard. Riders pay an exit fare based on time of day and distance traveled. Trips may include segments on multiple lines under one fare as long as the rider does not exit the faregates, with the exception of the "Farragut Crossing" out-of-station interchange between the and stations. At Farragut Crossing, riders may exit from one station and reenter at the other within 30 minutes on a single fare. When making a trip that uses Metrobus and Metrorail, a $2.25 discount is available when using a SmarTrip card when transferring from Metrobus to Metrorail, and Transfers from Metrorail to Metrobus are free; Transfers must be done within 2 hours.
When entering and exiting at the same station, users are normally charged a minimum fare ($2.25). However, since July 1, 2016, users have had a 15-minute grace period to exit the station; those who do so will receive a rebate of the amount paid as an autoload to their SmarTrip card.
Students at District of Columbia schools (public, charter, private, and parochial) ride both Metrobus and Metrorail for free.
Fare history
The contract for Metro's fare collection system was awarded in 1975 to Cubic Transportation Systems. Electronic fare collection using paper magnetic stripe cards started on July 1, 1977, a little more than a year after the first stations opened. Prior to electronic fare collection, exact change fareboxes were used. Metro's historic paper farecard system is also shared by Bay Area Rapid Transit, which Cubic won a contract for in 1974. Any remaining value stored on the paper cards was printed on the card at each exit, and passes were printed with the expiration date.
Several adjustments were made to shift the availability of passes from paper tickets to SmarTrip cards in 2012 and 2013. In May 2014 Metro announced plans to retrofit more than 500 fare vending machines throughout the system to dispense SmarTrip cards, rather than paper fare cards, and eventually eliminate magnetic fare cards entirely. This was completed in early December 2015 when the last paper farecard was sold. The faregates stopped accepting paper farecards on March 6, 2016, and the last day for trading in farecards to transfer the value to SmarTrip was June 30, 2016.
Safety and security
Security
Metro planners designed the system with passenger safety and order maintenance as primary considerations. The open vaulted ceiling design of stations and the limited obstructions on platforms allow few opportunities to conceal criminal activity. Station platforms are built away from station walls to limit vandalism and provide for diffused lighting of the station from recessed lights. Metro's attempts to reduce crime, combined with how the station environments were designed with crime prevention in mind, have contributed to Metro being among the safest and cleanest subway systems in the United States. There are nearly 6,000 video surveillance cameras used across the system to enhance security.
Metro is patrolled by its own police force, which is charged with ensuring the safety of passengers and employees. Transit Police officers patrol the Metro and Metrobus systems, and they have jurisdiction and arrest powers throughout the Metro service area for crimes that occur on or against transit authority facilities, or within of a Metrobus stop. The Metro Transit Police Department is one of two U.S. police agencies that has local police authority in three "state"-level jurisdictions (Maryland, Virginia, and the District of Columbia), the U.S. Park Police being the other.
Each city and county in the Metro service area has similar ordinances that regulate or prohibit vending on Metro-owned property, and which prohibit riders from eating, drinking, or smoking in Metro trains, buses, and stations; the Transit Police have a reputation for enforcing these laws rigorously. One widely publicized incident occurred in October 2000 when police arrested 12-year-old Ansche Hedgepeth for eating french fries in the station. In a 2004 opinion by John Roberts, now Chief Justice of the United States, the D.C. Circuit Court of Appeals upheld Hedgepeth's arrest. By then WMATA had answered negative publicity by adopting a policy of first issuing warnings to juveniles, and arresting them only after three violations within a year.
Metro's zero tolerance policy on food, trash and other sources of disorder embodies the "broken windows" philosophy of crime reduction. This philosophy also extends to the use of station restroom facilities. A longstanding policy, intended to curb unlawful and unwanted activity, has been to only allow employees to use Metro restrooms. One widely publicized example of this was when a pregnant woman was denied access to the bathroom by a station manager at the station. Metro now allows the use of restrooms by passengers who gain a station manager's permission, except during periods of heightened terror alerts.
On January 22, 2019, the D.C. Council voted 11–2 to override Mayor Muriel Bowser's veto of the Fare Evasion Decriminalization Act, setting the penalty for fare evasion at a $50 civil fine, a reduction from the previous criminal penalty of a fine up to $300 and 10 days in jail.
Random bag searches
On October 27, 2008, the Metro Transit Police Department announced plans to immediately begin random searches of backpacks, purses, and other bags. Transit police would search riders at random before boarding a bus or entering a station. It also explained its intent to stop anyone acting suspiciously. Metro claims that "Legal authority to inspect packages brought into the Metro system has been established by the court system on similar types of inspections in mass transit properties, airports, military facilities and courthouses." Metro Transit Police Chief Michael Taborn stated that, if someone were to turn around and simply enter the system through another escalator or elevator, Metro has "a plan to address suspicious behavior". Security expert Bruce Schneier characterized the plan as "security theater against a movie plot threat" and does not believe random bag searches actually improve security.
The Metro Riders' Advisory Council recommended to WMATA's board of directors that Metro hold at least one public meeting regarding the search program. , Metro had not conducted a single bag search.
In 2010 Metro once again announced that it would implement random bag searches, and conducted the first such searches on December 21, 2010. The searches consist of swabbing bags and packages for explosive residue, and X-raying or opening any packages which turned up positive. On the first day of searches, at least one false positive for explosives was produced, which Metro officials indicated could occur for a variety of reasons including if a passenger had recently been in contact with firearms or been to a firing range. The D.C. Bill of Rights Coalition and the Montgomery County Civil Rights Coalition circulated a petition against random bag searches, taking the position that the practice violates the Fourth Amendment to the United States Constitution and would not improve security. On January 3, 2011, Metro held a public forum for the searches at a Metro Riders' Advisory Council meeting, at which more than 50 riders spoke out, most of them in opposition to the searches. However at the meeting Metro officials called random bag inspections a "success" and claimed that few riders had complained.
After a prolonged absence, , bag searches have resumed at random stations throughout the Washington Metro area.
Safety
Accidents and incidents
Several collisions have occurred on Washington Metro, resulting in injuries and fatalities, along with numerous derailments with few or no injuries. WMATA has been criticized for disregarding safety warnings and advice from experts. The Tri-State Oversight Committee oversaw WMATA, but had no regulatory authority. Metro's safety department is usually in charge of investigating incidents, but could not require other Metro departments to implement its recommendations. Following several safety lapses, the Federal Transit Administration assumed oversight at WMATA.
Collisions
During the Blizzard of 1996, on January 6, a Metro operator was killed when a train failed to stop at the station. The four-car train overran the station platform and struck an unoccupied train that was awaiting assignment. The National Transportation Safety Board (NTSB) investigation found that the crash was a result of a failure in the train's computer-controlled braking system. The NTSB recommended that Metro grant train operators the ability to manually control the braking system, even in inclement weather, and recommended that Metro prohibit parked rail cars on tracks used by incoming outbound trains.
On November 3, 2004, an out-of-service Red Line train rolled backwards into the station, hitting an in-service train stopped at the platform. The rear car (1077) was telescoped by the first car of the standing train (4018). No one died, 20 people were injured. A 14-month investigation concluded that the train operator was most likely not alert as the train rolled backwards into the station. Safety officials estimated that had the train been full, at least 79 people would have died. The train operator was dismissed and Metro officials agreed to add rollback protection to more than 300 rail cars.
On June 22, 2009, at 5:02 pm, two trains on the Red Line collided. A southbound train heading toward Shady Grove stopped on the track short of the Fort Totten station and another southbound train collided with its rear. The front car of the moving train (1079) was telescoped by the rear car of the standing train (5066), and passengers were trapped. Nine people died and more than 70 were injured, dozens of whom were described as "walking wounded". Red Line service was suspended between the Fort Totten and Takoma stations, and New Hampshire Avenue was closed. One of the dead was the operator of the train that collided with the stopped train.
On November 29, 2009, at 4:27 am, two trains collided at the West Falls Church train yard. One train pulled in and collided with the back of the other train. No customers were aboard, and only minor injuries to the operators and cleaning staff were reported. However, three cars (1106, 1171, and 3216) were believed to be damaged beyond repair.
Derailments
On January 13, 1982, a train derailed at a malfunctioning crossover switch south of the station. In attempting to restore the train to the rails, supervisors failed to notice that another car had also derailed. The other rail car slid off the track and hit a tunnel support, killing three people and injuring 25 in its first fatal crash. Coincidentally, this crash occurred about 30 minutes after Air Florida Flight 90 crashed into the nearby 14th Street Bridge during a major snowstorm.
On January 20, 2003, during construction of a new canopy at the station, Metro began running trains through the center track even though it had not been constructed for standard operations, and a Blue Line train derailed at the switch. No injuries resulted but the crash delayed construction by a number of weeks.
On January 7, 2007, a Green Line train carrying approximately 120 people derailed near the station in downtown Washington. Trains were single-tracking at the time, and the derailment of the fifth car occurred where the train was switching from the south to northbound track. The crash injured at least 18 people and prompted the rescue of 60 people from a tunnel. At least one person had a serious but non-life-threatening injury. The incident was one of a series of five derailments involving 5000-series cars, with four of those occurring on side tracks and not involving passengers.
On June 9, 2008, an Orange Line train (2000-series) derailed between the Rosslyn and Court House stations.
On March 27, 2009, a Red Line train derailed just before 4:30 pm just south of station causing delays but no injuries. A second train was sent to move the first train but it too derailed when it was about from the first train.
On February 12, 2010, a Red Line train derailed at about 10:13 am as it left the station in downtown Washington. After leaving the station, the train entered the pocket track north of the station. As it continued, an automatic derailer at the end of the pocket track intentionally derailed the train as a safety measure. If the train had continued moving forward on the pocket track, it would have entered the path of an oncoming train. The wheels of the first two cars in the six-car, White-Flint-bound train were forced off the tracks, stopping the train. Almost all of the estimated 345 passengers were evacuated from the damaged train by 11:50 am and the NTSB arrived on the scene by noon. Two minor injuries were reported, and a third passenger was taken to George Washington University Hospital. The NTSB ruled the crash was due to the train operator's failure to follow standard procedures and WMATA management for failure to provide proper supervision of the train operator which resulted in the incomplete configuration of the train identification and destination codes leading to the routing of the train into the pocket track.
On April 24, 2012, around 7:15 pm, a Blue Line train bound for Franconia–Springfield derailed near Rosslyn. No injuries were reported.
On July 6, 2012, around 4:45 pm, a Green Line train bound for downtown Washington, D.C., and Branch Avenue derailed near West Hyattsville. No injuries were reported. A heat kink, due to the hot weather, was identified as the probable cause of the accident.
On August 6, 2015, a non-passenger train derailed outside the station. The track condition that caused the derailment had been detected a month earlier but was not repaired.
On July 29, 2016, a Silver Line train heading in the direction of Wiehle–Reston East station derailed outside East Falls Church station. Service was suspended between Ballston and West Falls Church and McLean stations on the Orange and Silver Lines.
On September 1, 2016, Metro announced the derailment of an empty six-car train in the Alexandria Rail Yard. No injuries or service interruptions were reported and an investigation is ongoing.
On January 15, 2018, a Red Line train derailed between the Farragut North and Metro Center stations. No injuries were reported. This was the first derailment of the new 7000-series trains.
On July 7, 2020, a 7000-series Red line train derailed one wheelset on departure from around 11:20 in the morning.
On October 12, 2021, a 7000-series Blue Line train derailed outside the Arlington Cemetery station. This forced the evacuation of all 187 passengers on board with no reported injuries. Cause of the derailment was initially stated to be an axle not up to specifications and resulted in sidelining the entire 7000-series fleet of trains, approximately 60% of WMATA's current trains through Friday, October 29, 2021, for further inspection. On October 28, 2021, WMATA announced that the system would continue running at a reduced capacity through November 15, 2021, as further investigation took place. The inspection determined a defect causes the car's wheels to be pushed outward. As of July 2022, the system was still running without most 7000-series cars. Workers manually inspect wheels on eight trains daily to catch the defect before it becomes problematic; the remaining cars are out of service pending an automated fix.
Safety measures
On July 13, 2009, WMATA adopted a "zero tolerance" policy for train or bus operators found to be texting or using other hand-held devices while on the job. This new and stricter policy came after investigations of several mass-transit accidents in the U.S. found that operators were texting at the time of the accident. The policy change was announced the day after a passenger of a Metro train videotaped the operator texting while operating the train.
Smoke incidents
During the early evening rush on January 12, 2015, a Yellow Line train stopped in the tunnel. It filled with smoke just after departing L'Enfant Plaza for Pentagon due to "an electrical arcing event" ahead in the tunnel. Everyone on board was evacuated; 84 people were taken to hospitals, and one died.
On March 14, 2016, an electrified rail caught fire between McPherson Square and Farragut West, causing significant disruptions on the Blue, Orange, and Silver lines. Two days later, the entire Metro system was shut down so its electric rail power grid could be inspected.
Future expansion
As of 2008, WMATA expects an average of one million riders daily by 2030. The need to increase capacity has renewed plans to add 220 cars to the system and reroute trains to alleviate congestion at the busiest stations. Population growth in the region has also revived efforts to extend service, build new stations, and construct additional lines.
Planned or proposed projects
Line extensions
The original plan called for ten future extensions on top of the core system. The Red Line would have been extended from the Rockville station northwest to Germantown, Maryland. The Green Line would have been lengthened northward from to Laurel, Maryland, and southward from to Brandywine, Maryland. The Blue Line initially consisted of a southwestern branch to Backlick Road and Burke, Virginia, which was never built. The Orange Line would have extended westward through Northern Virginia past the Vienna station to Centreville or Haymarket, and northeastward past to Bowie, Maryland. Alternatively, the Blue Line would have been extended east past to Bowie. The future Silver Line was also included in this proposal.
In 2001, officials considered realigning the Blue Line between and stations by building a bridge or tunnel from Virginia to a new station in Georgetown. Blue Line trains share a single tunnel with Orange Line and Silver Line trains to cross the Potomac River. The current tunnel limits service in each direction, creating a choke point. The proposal was later rejected due to cost, but Metro again started considering a similar scenario in 2011.
In 2005 the Department of Defense announced that it would be shifting 18,000 jobs to Fort Belvoir in Virginia and at least 5,000 jobs to Fort Meade in Maryland by 2012, as part of that year's Base Realignment and Closure plan. In anticipation of such a move, local officials and the military proposed extending the Blue and Green Lines to service each base. The proposed extension of the Green Line could cost $100 million per mile ($60 million per kilometer), and a light rail extension to Fort Belvoir was estimated to cost up to $800 million. Neither proposal has established timelines for planning or construction.
The Virginia Department of Transportation (VDOT) announced on January 18, 2008, that it and the Virginia Department of Rail and Public Transportation (DRPT) had begun work on a draft environmental impact statement (EIS) for the I-66 corridor in Fairfax and Prince William counties. According to VDOT the EIS, officially named the I-66 Multimodal Transportation and Environment Study, would focus on improving mobility along I-66 from the Capital Beltway (I-495) interchange in Fairfax County to the interchange with U.S. Route 15 in Prince William County. The EIS also allegedly includes a four-station extension of the Orange Line past Vienna. The extension would continue to run in the I-66 median and would have stations at Chain Bridge Road, Fair Oaks, Stringfellow Road and Centreville near Virginia Route 28 and U.S. Route 29. In its final report published June 8, 2012, the study and analysis revealed that an "extension would have a minimal impact on Metrorail ridership and volumes on study area roadways inside the Beltway and would therefore not relieve congestion in the study corridor."
In 2011 Metro began studying the needs of the system through 2040. WMATA subsequently published a study on the alternatives, none of which were funded for planning or construction. New Metro rail lines and extensions under consideration as part of this long-term plan included:
a new Loop line which parallels the Capital Beltway, known as the "Beltway Line"
a new Brown Line from the station to White Oak, Maryland, which would pass through the District and , running parallel to the Red Line.
rerouting the Yellow Line to either a new alignment, or a new tunnel parallel to the Green Line, in the District north of the Potomac River
a 5-station spur of the Green Line to National Harbor in Maryland
re-routing the Blue or Silver Lines in the District and/or building a separate express route for the Silver Line in Virginia
extensions to existing lines, including:
Red Line northwest to Metropolitan Grove (2 stations)
Orange Line east to Bowie (3 stations) or west to Centreville or Gainesville (3 or 5 stations, respectively)
Yellow Line south to Lorton (8 stations)
Green Line northeast to BWI Airport (6 stations) or southeast to White Plains (6 stations)
Blue Line east to Bowie (5 stations) or southwest to Potomac Mills (4 stations)
Silver Line northwest to Leesburg (3 stations)
four inter-line connections to allow greater service flexibility
several infill stations on existing lines
In September 2021, a report on the capacity improvements of Blue/Orange/Silver lines proposed four alternative extensions for the system:
Converting the Blue Line into a circle line, extending it to National Harbor and Alexandria. The proposed extension starts from a new station at , continues to Georgetown through a new tunnel under the Potomac River, then runs under M Street NE, just north of the existing Blue/Orange/Silver central segment, and connects to the Red Line at . It then turns south towards Buzzard Point, Joint Base Anacostia–Bolling, and National Harbor and crosses the Woodrow Wilson Bridge to Alexandria. The loop rejoins the current system at on the current Yellow Line, which is re-routed to .
A Blue Line extension to , which would follow a similar route through Georgetown to Union Station, then turn north towards Union Market and Ivy City before connecting with the Green Line at Greenbelt.
A Silver Line Express service from to with a similar route as the previous alternative.
A Silver Line extension to .
All four alternatives use the same central segment layout from Rosslyn to Union Station through Georgetown. NBC4 Washington further reported on the proposed loop in December 2022. At the time, there was a crowding problem at the Rosslyn station, and this expansion could be the solution to solve this crowding problem. A final design was published in July 2023.
Individual and infill stations
Before construction on Metro began, a proposed station was put forward for the Kennedy Center. Congress had already approved the construction of a station on the Orange/Blue/Silver Lines at 23rd and H Streets, near George Washington University, at the site of what is now Foggy Bottom station. According to a Washington Post article from February 1966, rerouting the line to accommodate a station under the center would cost an estimated $12.3 million. The National Capital Transportation Agency's administrator at the time, Walter J. McCarter, suggested that the Center "may wish to enhance the relationship to the station by constructing a pleasant, above-ground walkway from the station to the Center," referring to the then soon-to-be-built Foggy Bottom station. Rep. William B. Widnall, Republican of New Jersey, used it as an opportunity to push for moving the center to a central, downtown location.
The 2011 Metro transit-needs study identified five additional sites where infill stations could be built. These included Kansas Avenue and Montgomery College on the Red Line, respectively in Northwest D.C. and Rockville, Maryland; Oklahoma Avenue on the Blue, Orange, and Silver Lines near the D.C. Armory in Northeast D.C.; Eisenhower Valley on the Blue Line in Alexandria, Virginia; and the St. Elizabeths Hospital campus on the Green Line in Southeast D.C.
Related non-WMATA projects
A number of light rail and urban streetcar projects are under construction or have been proposed to extend or supplement service provided by Metro. The Purple Line, a light rail system, operated by the Maryland Transit Administration, is under construction as of 2024 and is scheduled to open in late 2027. The project was originally envisioned as a circular heavy rail line connecting the outer stations on each branch of the Metrorail system, in a pattern roughly mirroring the Capital Beltway. The current project will run between the and stations by way of and College Park. The Purple Line will connect both branches of the Red Line to the Green and Orange Lines, and would decrease the travel time between suburban Metro stations.
The Corridor Cities Transitway (CCT) is a proposed bus rapid transit line that would link Clarksburg, Maryland, in northern Montgomery County with the station on the Red Line. Assuming that the anticipated federal, state, and local government funds are provided, construction of the first of the system would begin in 2018.
In 2005, a Maryland lawmaker proposed a light rail system to connect areas of Southern Maryland, especially the rapidly growing area around the town of Waldorf, to the station on the Green Line.
The District of Columbia Department of Transportation is building the new DC Streetcar system to improve transit connectivity within the District. A tram line to connect Bolling Air Force Base to the station and was originally expected to open in 2010. Streetcar routes have been proposed in the Atlas District, Capitol Hill, and the K Street corridor. After seven years of construction, the Atlas District route, known as the H/Benning Street route, opened on February 27, 2016.
In 2013, the Georgetown Business Improvement District proposed a gondola lift between Georgetown and Rosslyn as an alternative to placing a Metro stop at Georgetown in its 2013–2028 economic plans. Washington, D.C., and Arlington County have been conducting feasibility studies for it since 2016.
In media
The Washington Metro has often appeared in movies and television shows set in Washington. However, due to fees and expenses required to film in the Metro, scenes of the Metro in film are often not of the Metro itself, but of other stand-in subway stations that are made to represent the Metro.
The vaulted ceilings of the Metro have become a cultural signifier of Washington, D.C., and are often seen in photographs and other art depicting the city.
| Technology | United States | null |
146181 | https://en.wikipedia.org/wiki/Tethys%20Ocean | Tethys Ocean | The Tethys Ocean ( ; ), also called the Tethys Sea or the Neo-Tethys, was a prehistoric ocean during much of the Mesozoic Era and early-mid Cenozoic Era. It was the predecessor to the modern Indian Ocean, the Mediterranean Sea, and the Eurasian inland marine basins (primarily represented today by the Black Sea and Caspian Sea).
During the early Mesozoic, as Pangaea broke up, the Tethys Ocean was defined as the ocean located between the ancient continents of Gondwana and Laurasia. After the opening of the Indian and Atlantic oceans during the Cretaceous Period and the breakup of these continents over the same period, it came to be defined as the ocean bordered by the continents of Africa, Eurasia, India, and Australasia. During the early-mid Cenozoic, the Indian, African, Australian and Arabian plates moved north and collided with the Eurasian plate, which created new borders to the ocean, a land barrier to the flow of currents between the Indian and Mediterranean basins, and the orogenies of the Alpide belt (including the Alps, Himalayas, Zagros, and Caucasus Mountains). All of these geological events, in addition to a drop in sea level rise from Antarctic glaciation, brought an end to the Tethys as it previously existed, fragmenting it into the Indian Ocean, the Mediterranean Sea, and the Paratethys.
It was preceded by the Paleo-Tethys Ocean, which lasted between the Cambrian and the Early Triassic, while the Neotethys formed during the Late Triassic and lasted in some form up to the Oligocene–Miocene boundary (about 24–21 million years ago) when it completely closed. A portion known as the Paratethys was isolated during the Oligocene (34 million years ago) and lasted up to the Pliocene (about 5 million years ago), when it largely dried out. The modern inland seas of Europe and Western Asia, namely the Black Sea and Caspian Sea, are remnants of the Paratethys Sea.
Etymology
The sea is named after Tethys, who, in ancient Greek mythology, is a water goddess, a sister and consort of Oceanus, mother of the Oceanid sea nymphs and of the world's great rivers, lakes and fountains.
Terminology and subdivisions
The eastern part of the Tethys Ocean is sometimes referred to as Eastern Tethys. The western part of the Tethys Ocean is called Tethys Sea, Western Tethys Ocean, or Paratethys or Alpine Tethys Ocean. The Black, Caspian, and Aral seas are thought to be its crustal remains, though the Black Sea may, in fact, be a remnant of the older Paleo-Tethys Ocean.
The Western Tethys was not simply a single open ocean. It covered many small plates, Cretaceous island arcs, and microcontinents. Many small oceanic basins (Valais Ocean, Piemont-Liguria Ocean, Meliata Ocean) were separated from each other by continental terranes on the Alboran, Iberian, and Apulian plates. The high sea level in the Mesozoic flooded most of these continental domains, forming shallow seas.
During the early Cenozoic, the Tethys Ocean could be divided into three sections: the Mediterranean Tethys (the direct predecessor to the Mediterranean Sea), the Peri-Tethys (a vast inland sea that covered much of eastern Europe and central Asia, and the direct predecessor to the Paratethys Sea), and the Indian Tethys (the direct predecessor to the Indian Ocean). The Turgai Strait extended out of the Peri-Tethys, connecting the Tethys with the Arctic Ocean.
As theories have improved, scientists have extended the "Tethys" name to refer to three similar oceans that preceded it, separating the continental terranes: in Asia, the Paleo-Tethys (Devonian–Triassic), Meso-Tethys (late Early Permian–Late Cretaceous), and Ceno-Tethys (Late-Triassic–Cenozoic) are recognized. None of the Tethys oceans should be confused with the Rheic Ocean, which existed to the west of them in the Silurian Period. To the north of the Tethys, the then-land mass is called Angaraland and to the south of it, it is called Gondwanaland.
Modern theory
From the Ediacaran (600 ) into the Devonian (360 ), the Proto-Tethys Ocean existed and was situated between Baltica and Laurentia to the north and Gondwana to the south.
From the Silurian (440 ) through the Jurassic periods, the Paleo-Tethys Ocean existed between the Hunic terranes and Gondwana. Over a period of 400 million years, continental terranes intermittently separated from Gondwana in the Southern Hemisphere to migrate northward to form Asia in the Northern Hemisphere.
Triassic Period
About 250 Mya, during the Triassic, a new ocean began forming in the southern end of the Paleo-Tethys Ocean. A rift formed along the northern continental shelf of Southern Pangaea (Gondwana). Over the next 60 million years, that piece of shelf, known as Cimmeria, traveled north, pushing the floor of the Paleo-Tethys Ocean under the eastern end of northern Pangaea (early / proto- Laurasia). The Neo-Tethys Ocean formed between Cimmeria and Gondwana, directly over where the Paleo-Tethys formerly rested.
Jurassic Period
During the Jurassic period about 150 Mya, Cimmeria finally collided with Laurasia and stalled, so the ocean floor behind it buckled under, forming the Tethys Trench. Water levels rose, and the western Tethys shallowly covered significant portions of Europe, forming the first Tethys Sea. Around the same time, Laurasia and Gondwana began drifting apart, opening an extension of the Tethys Sea between them which today is the part of the Atlantic Ocean between the Mediterranean and the Caribbean. As North and South America were still attached to the rest of Laurasia and Gondwana, respectively, the Tethys Ocean in its widest extension was part of a continuous oceanic belt running around the Earth between about latitude 30°N and the Equator. Thus, ocean currents at the time around the Early Cretaceous ran very differently from the way they do today.
Late Cretaceous
Between the Jurassic and the Late Cretaceous, which started about 100 Mya, Gondwana began breaking up, pushing Africa and India north across the Tethys and opening up the Indian Ocean.
Cenozoic
Throughout the Cenozoic (66 million to the dawn of the Neogene, 23 Mya), the connections between the Atlantic and Indian Oceans across the Tethys were eventually closed off in what is now the Middle East during the Miocene, as a consequence of the northern migration of Africa/Arabia and global sea levels falling due to the concurrent formation of the Antarctic Ice Sheet. This decoupling occurred in two steps, first around 20 Mya and another around 14 Mya. The complete closure of the Tethys led to a global reorganization of currents, and is what is thought to have allowed for upwelling in the Arabian Sea and led to the establishment of the modern South Asian Monsoon. It also caused major modifications to the functioning of the AMOC and ACC.
During the Oligocene (33.9 to 23 Mya), large parts of central and eastern Europe were covered by a northern branch of the Tethys Ocean, called the Paratethys. The Paratethys was separated from the Tethys with the formation of the Alps, Carpathians, Dinarides, Taurus, and Elburz mountains during the Alpine orogeny. During the late Miocene, the Paratethys gradually disappeared, and became an isolated inland sea. Separation from the wider Tethys during the early Miocene initially led to a boost in primary productivity for the Paratethys, but this gave way to a total ecosystem collapse during the late Miocene as a result of rapid dissolution of carbonate.
Historical theory
In Chapter 13 of his 1845 book, Roderick Murchison described a distinctive formation extending from the Black Sea to the Aral Sea in which the creatures differed from those of the purely marine period that preceded them. The Miocene deposits of Crimea and Taman (south of the Sea of Azov) are identical with formations surrounding the present Caspian Sea, in which the univalves of freshwater origin are associated with forms of Cardiacae and Mytili that are common to partially saline or brackish waters. This distinctive fauna has been found throughout all the enormously developed Tertiary formations of the southern and south-eastern steppes.
On the accompanying map, Murchison shows the Aralo-Caspian Formation extending from close to the Danube delta across Crimea, up the east side of the Volga river to Samara, then south of the Urals to beyond the Aral Sea. Brackish and upper freshwater components (OSM) of the Miocene are now known to extend through the North Alpine foreland basin and onto the Swabian Jura with thickness of up to ; these were deposited in the Paratethys when the Alpine front was still farther south.
In 1885, the Austrian palaeontologist Melchior Neumayr deduced the existence of the Tethys Ocean from Mesozoic marine sediments and their distribution, calling his concept () and described it as a Jurassic seaway, which extended from the Caribbean to the Himalayas.
In 1893, the Austrian geologist Eduard Suess proposed the hypothesis that an ancient and extinct inland sea had once existed between Laurasia and the continents which formed Gondwana II. He named it the Tethys Sea after the Greek sea goddess Tethys. He provided evidence for his theory using fossil records from the Alps and Africa. He proposed the concept of Tethys in his four-volume work (The Face of the Earth).
In the following decades during the 20th century, "mobilist" geologists such as Uhlig (1911), Diener (1925), and Daque (1926) regarded Tethys as a large trough between two supercontinents which lasted from the late Palaeozoic until continental fragments derived from Gondwana obliterated it.
After World War II, Tethys was described as a triangular ocean with a wide eastern end.
From 1920s to the 1960s, "fixist" geologists, however, regarded Tethys as a composite trough, which evolved through a series of orogenic cycles. They used the terms 'Paleotethys', 'Mesotethys', and 'Neotethys' for the Caledonian, Variscan, and Alpine orogenies, respectively. In the 1970s and 1980s, these terms and 'Proto-Tethys', were used in different senses by various authors, but the concept of a single ocean wedging into Pangea from the east, roughly where Suess first proposed it, remained.
In the 1960s, the theory of plate tectonics became established, and Suess's "sea" could clearly be seen to have been an ocean. Plate tectonics provided an explanation for the mechanism by which the former ocean disappeared: oceanic crust can subduct under continental crust.
Tethys was considered an oceanic plate by Smith (1971); Dewey, Pitman, Ryan and Bonnin (1973); Laubscher and Bernoulli (1973); and Bijou-Duval, Dercourt and Pichon (1977).
| Physical sciences | Paleogeography | Earth science |
146227 | https://en.wikipedia.org/wiki/Sports%20car | Sports car | A sports car is a type of car that is designed with an emphasis on dynamic performance, such as handling, acceleration, top speed, the thrill of driving, and racing capability. Sports cars originated in Europe in the early 1910s and are currently produced by many manufacturers around the world.
Definition
Definitions of sports cars often relate to how the car design is optimised for dynamic performance, without any specific minimum requirements; both a Triumph Spitfire and Ferrari 488 Pista can be considered sports cars, despite vastly different levels of performance. Broader definitions of sports cars include cars "in which performance takes precedence over carrying capacity", or that emphasise the "thrill of driving" or are marketed "using the excitement of speed and the glamour of the (race)track" However, other people have more specific definitions, such as "must be a two-seater or a 2+2 seater" or a car with two seats only.
In the United Kingdom, early recorded usage of the "sports car" was in The Times newspaper in 1919. The first known use of the term in the United States was in 1928. Sports cars started to become popular during the 1920s. The term initially described two-seat roadsters (cars without a fixed roof), however, since the 1970s the term has also been used for cars with a fixed roof (which were previously considered grand tourers).
Attributing the definition of 'sports car' to any particular model can be controversial or the subject of debate among enthusiasts. Authors and experts have often contributed their ideas to capture a definition. Insurance companies have also attempted to use mathematical formulae to categorise sports cars, often charging more for insurance due to the inherent risk of performance driving.
There is no fixed distinction between sports cars and other categories of performance cars, such as muscle cars and grand tourers, with some cars being members of several categories.
Common characteristics
Seating layout
Traditionally, the most common layout for sports cars was a roadster (a two-seat car without a fixed roof). However, there are also several examples of early sports cars with four seats.
Sports cars are not usually intended to transport more than two adult occupants regularly, so most modern sports cars are generally two-seat or 2+2 layout (two smaller rear seats for children or occasional adult use). Larger cars with more spacious rear-seat accommodation are usually considered sports sedans rather than sports cars.
The 1993–1998 McLaren F1 is notable for using a three-seat layout, where the front row consists of a centrally-located driver's seat.
Engine and drivetrain layout
The location of the engine and driven wheels significantly influence the handling characteristics of a car and are therefore crucial in the design of a sports car. Traditionally, most sports cars have used rear-wheel drive with the engine either located at the front (FR layout) or in the middle of the vehicle (MR layout). Examples of FR layout sports cars include the Caterham 7, Mazda MX-5, and the Dodge Viper. Examples of MR layout sports cars are the Ferrari 488, Ford GT, and Toyota MR2. To avoid a front-heavy weight distribution, many FR layout sports cars are designed so that the engine is located further back in the engine bay, as close to the firewall as possible.
Since the 1990s, all-wheel drive has become more common in sports cars. All-wheel drive offers better acceleration and favorable handling characteristics (especially in slippery conditions), but is often heavier and more mechanically complex than traditional layouts. Examples of all-wheel drive sports cars are the Lamborghini Huracan, Bugatti Veyron, and Nissan GT-R.
Rear engine layouts are not typical for sports cars, with the notable exception of the Porsche 911.
The front-wheel drive layout with the engine at the front (FF layout) is generally the most common for cars, but it is not as common among traditional sports cars. Nonetheless, the FF layout is used by sport compacts and hot hatches such as the Mazdaspeed3. Sports cars with an FF layout include the Fiat Barchetta, Saab Sonett, or Opel Tigra.
Europe
1895–1917: Brass Era of cars
The basis for the sports car is traced to the early 20th century touring cars and roadsters, and the term 'sports car' would not be coined until after World War One.
A car considered to be "a sports-car years ahead of its time" is the 1903 Mercedes Simplex 60 hp, described at the time as a fast touring car and designed by Wilhelm Maybach and Paul Daimler. The Mercedes included pioneering features such as a pressed-steel chassis, a gated 4-speed transmission, pushrod-actuated overhead inlet valves, a honeycomb radiator, low-tension magneto ignition, a long wheelbase, a low center of mass and a very effective suspension system. The overall result was a "safe and well-balanced machine" with a higher performance than any other contemporary production car. At the 1903 Gordon Bennett Cup, a production Simplex 60 hp was entered only due to a specially-built 90 hp racing car being destroyed in a fire; the 60 hp famously went on to win the race.
The 1910 Austro-Daimler 27/80 is another early sports car which had success in motor racing. The 27/80 was designed by Ferdinand Porsche, who drove the car to victory in the 1910 Prince Henry Tour motor race. The Vauxhall and Austro-Daimler —like the Mercedes Simplex 60 hp— were production fast touring cars. The 1912 Hispano-Suiza Alfonso XIII is also considered one of the earliest sports cars, as it was a "purpose built, high performance, two-seater production automobile". The model was named after King Alfonso XIII of Spain, a patron of the car's chief designer and an enthusiast for the marque. Other early sports cars include the 1905 Isotta Fraschini Tipo D, the 1906 Rolls-Royce Silver Ghost, the 1908 Delage, the 1910 Bugatti Type 13, and the 1912 DFP 12/15.
Early motor racing events included the 1903 Paris–Madrid race, the 1905–1907 Herkomer Trophy, the 1908-1911 Prince Henry Tour and the 1911–present Monte Carlo Rally. The Prince Henry Tours (which were similar to modern car rallies) were among the sporting events of the period, bringing renown to successful entrants. The Prince Henry Tours started the evolution of reasonably large and technically advanced production sports cars.
In England, the development of sporting cars was inhibited by the Motor Car Act 1903, which imposed a speed limit of on all public roads. This led to the 1907 opening of the Brooklands motor circuit, which inspired the development of performance cars such as the 1910 Vauxhall Prince Henry, 1910 Sunbeam 12/16, 1910 Talbot 25 hp, 1910 Straker-Squire 15 hp and 1913 Star 15.9 hp.
1919–1929: Vintage Era cars
Following the halt in sports car production caused by World War I, Europe returned to manufacturing automobiles from around 1920. It was around this time that the term 'Sports Car' began to appear in the motor catalogues, although the exact origin of the name is not known. The decade that followed became known as the vintage era and featured rapid technical advances over the preceding Brass Era cars. Engine performance benefited from the abandonment of "tax horsepower" (where vehicles were taxed based on bore and number of cylinders, rather than actual power output) and the introduction of leaded fuel, which increased power by allowing for higher compression ratios.
In the early 1920s, the cost to produce a racing car was not significantly higher than a road car, therefore several manufacturers used the design from the current year's racing car for the next year's sports car. For example, the 1921 Ballot 2LS based on the racing car that finished third at the 1921 French Grand Prix. The Benz 28/95PS was also a successful racing car, with victories including the 1921 Coppa Florio.
Another approach— such as that used by Morris Garages— was to convert touring cars into sports cars.
The first 24 Hours of Le Mans race for sports cars was held in 1923, although the two-seat sports cars only competed in the smallest class, with the majority of cars entered being four-seat fast touring cars. "This race, together with the Tourist Trophy Series of Races, organised after the first World War by the R.A.C., appealed to the public imagination and offered to the manufacturers of the more sporting cars an excellent opportunity for boosting sales of their products." The classic Italian road races— the Targa Florio, and the Mille Miglia (first held in 1927)— also captured the public's imagination.
By 1925, the higher profits available for four-seater cars resulted in the production of two-seat sports cars being limited to smaller manufacturers such as Aston-Martin (350 Astons built from 1921 to 1939) and Frazer-Nash (323 cars built from 1924 to 1939). Then by the late 1920s, the cost of producing racing cars (especially Grand Prix cars) escalated, causing more manufacturers to produce cars for the growing sports car market instead.
Significant manufacturers of sports cars in the late 1920s were AC Cars, Alfa Romeo, Alvis, Amilcar, Bignan and Samson, Chenard-Walcker, Delage, Hispano-Suiza, Hotchkiss, Mercedes-Benz and Nazzaro. Two cars from the Vintage Era that would influence sports cars for many years were the Austin Seven and MG M-type "Midget". Successful sports cars from Bentley during this era were the Bentley 3 Litre (1921–1929) and the Bentley Speed Six (1928–1930), with the former famously described by Bugatti's founder as "the fastest lorry in the world".
1930–1939: Pre-war Era cars
Between the Great Depression and the World War II the pre-war era was a period of decline in importance for sports car manufacturers, although the period was not devoid of advances, for example streamlining. Cheap, light-weight family sedans with independent front suspension— such as the BMW 303, Citroën Traction Avant and Fiat 508— offered similar handling and comfort to the more expensive sports cars. Powerful, reliable, and economical (although softly suspended) American saloons began to be imported to Europe in significant numbers. Sports car ownership was increased through models such as the Austin 7 and Wolseley Hornet six, however many of these sports cars did not offer any performance upgrades over the mass-produced cars upon which they were based.
The highest selling sports car company of the 1930s was Morris Garages, who produced 'MG Midget' models of the M-Type, J-Type, P-Type and T-Type. The K3 version of the K-Type Magnette was a successful racing car, achieving success in the Mille Miglia, Tourist Trophy and 24 Hours of Le Mans.
The Bugatti Type 57 (1934-1940) was another significant sports car of the pre-war era and is now among the most valuable cars in the world. The T57 was successful in sports car races, including winning the 1937 24 Hours of Le Mans and 1939 24 Hours of Le Mans. Another successful Bugatti sports car was the Bugatti Type 55 (1932-1935), which was based on the Type 51 Grand Prix racing car.
1939–1959: Expansion following World War II
The decade following the Second World War saw an "immense growth of interest in the sports car, but also the most important and diverse technical developments [and] very rapid and genuine improvement in the qualities of every modern production car; assisted by new design and manufacturing techniques a consistently higher level of handling properties has been achieved."
In Italy, a small but wealthy market segment allowed for the manufacture of a limited number of high-performance models directly allied to contemporary Grand Prix machines, such as the 1948 Ferrari 166 S. A new concept altogether was the modern Gran Turismo class from Italy, which was in effect unknown before the war: sustained high-speed motoring from relatively modest engine size and compact closed or berlinetta coachwork. The 1947 Maserati A6 1500 two-seat berlinetta was the first production model from Maserati.
In Germany, the motor industry was devastated by the war, but a small number of manufacturers returned it to prominence. In 1948, the Porsche 356 was released as the debut model from Porsche. The significance of the Porsche 356 and its successors was described in 1957 as "future historians must see them as among the most important of mid-century production cars". The 1954 Mercedes-Benz 300 SL is another significant car from this era.
1960–1979: Lightweight roadsters, mid-engined supercars
The 1961 Jaguar E-Type is an iconic sports car of the early 1960s, due to its attractive styling and claimed top speed of . The E-type was produced for 14 years and was initially powered by a six-cylinder engine, followed by a V12 engine for the final generation.
In 1962, the MG B introduced a new era of affordable lightweight four-cylinder roadsters. The MG B used a unibody construction and was produced until 1980. Other successful lightweight roadsters include the Triumph Spitfire (1962-1980) and the Alfa Romeo Spider (1966-1993). The Fiat X1/9 (1972-1989) was unusual for its use of a mid-engine design in an affordable roadster model. A late entrant to the affordable roadster market was the 1975 Triumph TR7, however by the late 1970s the demand for this style of car was in decline, resulting in production ceasing in 1982.
The original Lotus Elan (1962-1975) two-seat coupe and roadster models are an early commercial success for the philosophy of achieving performance through minimizing weight and has been rated as one of the top 10 sports cars of the 1960s. The Elan featured fibreglass bodies, a backbone chassis, and overhead camshaft engines.
A different style of roadster was the AC Cobra, released in 1962, which was fitted with V8 engines up to in size by Shelby.
The Porsche 911 was released in 1964 and has remained in production since. The 911 is notable for its use of the uncommon rear-engine design and the use of a flat-six engine. Another successful rear-engine sports car was the original Alpine A110 (1961-1977), which was a successful rally car during the Group 4 era.
In 1965, the BMW New Class Coupes were released, leading to the BMW 6 Series which remains in production to this day.
The Lamborghini Miura (1966) and Alfa Romeo 33 Stradale (1967) mid-engined high-performance cars are often cited as the first supercars. Other significant European models of the 1960s and 1970s which might be considered supercars today are the Ferrari 250 GTO (1962-1964), Ferrari 250 GT Lusso (1963-1964), Ferrari 275 GTB/4 (1966-1968), Maserati Ghibli (1967-1973), Ferrari Daytona (1968-1973), Dino 246 (1969-1974), De Tomaso Pantera (1971-1993), Ferrari 308 GTB (1975-1980) and BMW M1 (1978-1981).
In 1966, the Jensen FF became the first sports car to use all-wheel drive.
The Ford Capri is a 2+2 coupe that was produced from 1968 to 1986 and intended to be a smaller European equivalent of the Ford Mustang. A main rival to the Capri was Opel Manta, which was produced from 1970 to 1988.
The 1973-1978 Lancia Stratos was a mid-engined two-seat coupe that was powered by a Ferrari V6 engine. This was an unusual arrangement for a car used to compete in rallying, nonetheless it was very successful and won the World Rally Championship in 1974, 1975, and 1976.
The Lancia Montecarlo was produced from 1975 to 1981 and is a mid-engine two-seater, available as a coupé or a targa-top. It was sold as Lancia Scorpion in the USA. Its racing variant, Montecarlo Turbo, won the 1979 World Championship for Makes in its division and overall for 1980 World Championship for Makes and 1981 World Endurance Championship for Makes. Montecarlo also won the 1980 Deutsche Rennsport Meisterschaft and Giro d'Italia automobilistico marathon. The Montecarlo was a basis for the silhouette racing car, Lancia Rally 037.
In the 1970s, turbocharging began to be adopted by sports cars, such as the BMW 2002 Turbo in 1973, the first Porsche 911 Turbo in 1975, and the Saab 99 Turbo in 1978.
1980–1999: Turbocharging and all-wheel drive emerge
Turbocharging became increasingly popular in the 1980s, from relatively affordable coupes such as the 1980–1986 Renault Fuego and 1992–1996 Rover 220 Coupé Turbo, to expensive supercars such as the 1984-1987 Ferrari 288 GTO and 1987-1992 Ferrari F40.
In the late 1980s and early 1990s, several manufacturers developed supercars that competed for production car top speed records. These cars included the 1986–1993 Porsche 959, 1991–1995 Bugatti EB 110, 1992–1994 Jaguar XJ220 and 1993–98 McLaren F1.
The 1980-1995 Audi Quattro was a pioneering all-wheel drive sports car. The 1995 Porsche 911 Turbo (993) saw the 911 Turbo model switch to all-wheel drive, a drivetrain layout that the model uses to this day.
The BMW M3 was released in 1986 and has been produced for every generation since. The 1993-1996 Mercedes-Benz W124 E36 AMG was the mass-produced AMG model. Audi's equivalent division, called "RS", was launched in 1994 with the Audi RS 2 Avant.
Ford Europe withdrew from the sports car market at the end of 1986 when the Capri was discontinued after a production run of nearly two decades. There was no direct successor, as Ford was concentrating on higher-performance versions of its hatchback and saloon models at the time.
In 1989, a new generation of Lotus Elan roadster was released which used a front-wheel drive layout, a controversial choice for a "purist" sports car. The Elan sold poorly and was discontinued after three years. The 1996 Lotus Elise, a mid-engined, rear-wheel drive roadster, was much more successful and remained in production until 2021.
Roadsters enjoyed a resurgence in the mid-1990s, including the 1989-present Mazda MX-5, the 1995-2002 BMW Z3 (succeeded by the 2002-2016 BMW Z4), the 1995-2002 MG F, the 1996–present Porsche Boxster and the 1998–present Audi TT.
The Honda S2000 roadster was introduced in 1999 for the 2000 model year and was noted for its high-revving 4-cylinder engine and its exceptionally high specific output of 125 horsepower per litre.
2000–present: Turbos become dominant, hybrids emerge
The 2000–2021 Lotus Exige was introduced as a coupe version of Elise. Similarly, Porsche Cayman (987) was introduced in 2006 as the coupe equivalent to the Porsche Boxster roadster. Lotus also expanded its model range with the 2009–2021 Lotus Evora, a larger four-seat coupe.
Audi's first mid-engined supercar is the 2006–present Audi R8. Other sports cars of the 2000s were the 2005-2010 Alfa Romeo Brera/Spider, 2009-2015 Peugeot RCZ, and the 2008-2017 reintroduction of the Volkswagen Scirocco (a coupe based on the VW Golf platform).
Reflecting overall car industry trends, the mid-2010s saw naturally aspirated engines being replaced by turbocharged engines. Ferrari's first regular production turbocharged engine was used in the 2014-2017 Ferrari California T, followed by the 2015-2019 Ferrari 488. Similarly, in 2016, the Porsche 911 (991.2) began to use turbocharging on all models and the Porsche 982 Cayman/Boxster downsized from a six-cylinder engine to a turbocharged four-cylinder engine.
Also in the 2010s, dual-clutch transmissions became more widespread, causing manual transmissions to decline in sales and no longer be offered on some models.
Hybrid-electric sports cars began to appear in the 2010s— notably the 2013-2016 LaFerrari, 2013-2015 McLaren P1, 2013-2015 Porsche 918 Spyder "hypercars". The 2014–2020 BMW i8 was also an early plug-in hybrid sports car.
McLaren began permanent car manufacturing operations with the 2011-2014 McLaren 12C.
In 2013, the Jaguar F-Type saw the brand return to the two-seat sports car market, with the four-seat grand tourer Jaguar XK discontinued the following year.
The BMW 2 Series coupe and convertible were introduced in 2013 to sit below the larger BMW 4 Series models, with the new BMW M2 high-performance model introduced in 2015.
The 2013–present Alfa Romeo 4C two-seat coupe and roadster used a carbon-fibre body and became Alfa's first mid-engine sports car since the 1967-1969 Alfa Romeo 33 Stradale.
Fiat had exited the roadster market with the end of Fiat Barchetta production in 2005. The company resumed production of roadsters in 2016 with the Fiat 124 Spider, which is based on the Mazda MX-5.
In 2017, Renault revived the Alpine brand for the 2017–present Alpine A110 mid-engine coupe.
United States
During the 1910s and 1920s, American manufacturers of smaller sports cars included Apperson, Kissel, Marion, Midland, National, Overland, Stoddard-Dayton and Thomas; manufacturers of larger sports cars included Chadwick, Mercer, Stutz Motor Company, and Simplex.
Since the 1960s, American performance cars have often been designed as muscle cars rather than sports cars. However, several American two-seat sports cars have also been produced, such as the 1953–present Chevrolet Corvette, 1962-1967 Shelby Cobra, 1983-1988 Pontiac Fiero, 1991-2017 Dodge Viper, and 2005-2006 Ford GT.
Asia
1959–1968: Beginnings
The first Japanese sports car was the 1959-1960 Datsun 211, a two-seat roadster built on the chassis of a compact pickup truck and powered by a engine. Only 20 cars were built, and the 1963-1965 Datsun SP310— based on the chassis of a passenger sedan instead of a pickup truck— is often considered Datsun's first mass-production sports car.
Honda's first sports car was the 1963-1964 Honda S500, a two-seat roadster with independent suspension for all wheels and a DOHC engine. In 1965, Toyota joined the two-seat roadster market with the Toyota Sports 800.
Mazda is noted for its use of rotary engines, beginning in 1967 with the Mazda Cosmo. The Cosmo was a two-seat coupe with a rotary engine producing up to . Mazda continued to produce sports cars with rotary engines (sometimes turbocharged) until the Mazda RX-8 ended production in 2012.
The Toyota 2000GT, produced from 1967 to 1970, was an expensive two-seat coupe that greatly changed overseas perceptions of the Japanese automotive industry. The 2000GT demonstrated that Japan was capable of producing high-end sports cars to rival the traditional European brands.
1969–1977: Mass production begins
In 1969, Nissan introduced the Nissan Fairlady Z / Datsun 240Z two-seat coupe, powered by a six-cylinder engine and described as providing similar performance to the Jaguar E-Type at a more affordable price. The 240Z began the lineage of Nissan "Z cars" which continues through to today's Nissan Z (RZ34). In 1974, Nissan expanded their coupe range with the Nissan Silvia 2+2 coupe, which was powered by a four-cylinder engine and produced until 2002.
Also in 1969, Mitsubishi's first performance car was introduced, in the form of the Mitsubishi Colt 11-F Super Sports coupe. The 11-F Super Sports was followed by the 1970-1977 Mitsubishi Galant GTO and 1971-1975 Mitsubishi Galant FTO, both based on a platform shared with the Galant sedan.
Toyota's mass-production 2+2 coupes of the 1970s consisted of the Celica, Supra, Corolla Levin, and Sprinter Trueno. The Celica was introduced in 1971 and remained in production until 2006. From 1979 to 1986, the Supra name was used for six-cylinder versions of the Celica, until the Supra moved to a separate platform from 1986 to 2002. The Corolla Levin / Sprinter Trueno were based on the Toyota Corolla hatchback platform and produced from 1972 to 2000.
The Nissan Skyline GT-R was initially produced as a sedan for two years before a coupe model was introduced in 1971. This first generation Skyline GT-R had rear-wheel drive, a six-cylinder engine, and was produced until 1972.
1978–1988: Front-wheel drive introduced
The Honda Prelude front-wheel drive 2+2 coupe was launched in 1978 and remained in production until 2001. The 1985-2006 Honda Integra was also a front-wheel drive 2+2 coupe produced by Honda. Other 2+2 models included the 1982-1989 Mitsubishi Starion (turbocharged and rear-wheel drive) and the 1985-1991 Subaru XT (available with a turbocharger and all-wheel drive). Subaru has produced few sports cars in its history, instead focusing on rally-influenced sedans/hatchbacks for their performance models, such as the Liberty RS and Impreza WRX/STi models.
In 1984, the Toyota MR2 two-seat coupe became Japan's first production mid-engine car. The MR2 switched to a two-seat roadster body style for the final generation from 1999 to 2007.
The first Korean coupe model was the 1988 Hyundai Scoupe, which used front-wheel drive and was based on the Excel hatchback. The Scoupe was followed by 1996-2008 Hyundai Tiburon and 2011-2022 Hyundai Veloster.
1989–2006: All-wheel drive, first supercars
In the 1990s, multiple Japanese automakers made flagship sports cars, such as the Toyota Supra, Nissan Skyline GT-R, Honda NSX, Mazda RX-7, and Mitsubishi 3000GT, which notably performed well against their European competition. These automakers had a well-documented gentlemen's agreement to officially limit advertised power figures for these vehicles to a maximum of .
The Nissan Skyline GT-R was reintroduced in 1989-2002 (R32, R33, and R34 generations) which became famous for their use of turbocharging and all-wheel drive, which provided performance comparable with many more expensive sports cars.
The 1990-2005 Honda NSX is considered Japan's first supercar. The NSX was praised for being more reliable and user-friendly than contemporary European supercars.
The Honda S2000 is an open top sports car that was manufactured from 1999 to 2009. The S2000 is named for its engine displacement of two liters, carrying on in the tradition of the S500, S600, and S800 roadsters of the 1960s. Its engine is notable for its high specific power output.
The Mitsubishi GTO coupe/convertible was introduced in 1990. The base models used front-wheel drive and a naturally aspirated V6 engine, however all-wheel drive and a turbocharged V6 engine were also available. To sit below the GTO in the model range, the Mitsubishi FTO front-wheel drive coupe was introduced in 1994. Both the GTO and FTO were discontinued in 2000.
Suzuki's first sports car was the 1991-1998 Suzuki Cappuccino, a two-seat roadster kei car with rear-wheel drive and a turbocharged engine.
From 2003 to 2012 Mazda produced the Mazda RX-8, a rear-wheel drive quad coupé powered by a 1.3 L Renesis twin-rotor engine.
2007–present: Declining popularity of coupes
Due to production constraints, lower demand, and environmental regulations, the viability of new Japanese sports cars began to decrease in the mid-2000s. The latest generation (R35) of the Nissan Skyline GT-R started production in 2007 as the Nissan GT-R. It was noteworthy for offering supercar performance with sports car practicality.
The Lexus LFA supercar was released by Lexus in 2010, a two-seat front-engine coupe powered by a V10 engine.
The Toyota 86 / Subaru BRZ is a 2+2 coupe that was introduced in 2012 and currently remains in production with a new model released for the 2022 model year. The 86/BRZ is a rare modern example of a relatively affordable rear-wheel drive sports car.
The 2016–2022 Honda NSX (2nd generation) supercar marked a change in approach for Honda, by using all-wheel drive, a hybrid drivetrain, turbocharging, and a dual-clutch transmission.
Toyota relaunched the Supra nameplate in 2019 after a 17-year hiatus for the Toyota GR Supra front-engine, rear-wheel-drive coupe. Nissan also released the new RZ34 generation of their Nissan Z in 2022.
| Technology | Motorized road transport | null |
146230 | https://en.wikipedia.org/wiki/Iceberg | Iceberg | An iceberg is a piece of fresh water ice more than long that has broken off a glacier or an ice shelf and is floating freely in open water. Smaller chunks of floating glacially derived ice are called "growlers" or "bergy bits". Much of an iceberg is below the water's surface, which led to the expression "tip of the iceberg" to illustrate a small part of a larger unseen issue. Icebergs are considered a serious maritime hazard.
Icebergs vary considerably in size and shape. Icebergs that calve from glaciers in Greenland are often irregularly shaped while Antarctic ice shelves often produce large tabular (table top) icebergs. The largest iceberg in recent history, named B-15, was measured at nearly in 2000. The largest iceberg on record was an Antarctic tabular iceberg measuring sighted west of Scott Island, in the South Pacific Ocean, by the USS Glacier on November 12, 1956. This iceberg was larger than Belgium.
Etymology
The word iceberg is a partial loan translation from the Dutch word ijsberg, literally meaning ice mountain, cognate to Danish isbjerg, German Eisberg, Low Saxon Iesbarg and Swedish isberg.
Overview
Typically about one-tenth of the volume of an iceberg is above water, which follows from Archimedes's Principle of buoyancy; the density of pure ice is about 920 kg/m3 (57 lb/cu ft), and that of seawater about . The contour of the underwater portion can be difficult to judge by looking at the portion above the surface.
The largest icebergs recorded have been calved, or broken off, from the Ross Ice Shelf of Antarctica. Icebergs may reach a height of more than above the sea surface and have mass ranging from about 100,000 tonnes up to more than 10 million tonnes. Icebergs or pieces of floating ice smaller than 5 meters above the sea surface are classified as "bergy bits"; smaller than 1 meter—"growlers". The largest known iceberg in the North Atlantic was above sea level, reported by the USCG icebreaker Eastwind in 1958, making it the height of a 55-story building. These icebergs originate from the glaciers of western Greenland and may have interior temperatures of .
Drift
A given iceberg's trajectory through the ocean can be modelled by integrating the equation
where m is the iceberg mass, v the drift velocity, and the variables f, k, and F correspond to the Coriolis force, the vertical unit vector, and a given force. The subscripts a, w, r, s, and p correspond to the air drag, water drag, wave radiation force, sea ice drag, and the horizontal pressure gradient force.
Icebergs deteriorate through melting and fracturing, which changes the mass m, as well as the surface area, volume, and stability of the iceberg. Iceberg deterioration and drift, therefore, are interconnected ie. iceberg thermodynamics, and fracturing must be considered when modelling iceberg drift.
Winds and currents may move icebergs close to coastlines, where they can become frozen into pack ice (one form of sea ice), or drift into shallow waters, where they can come into contact with the seabed, a phenomenon called seabed gouging.
Mass loss
Icebergs lose mass due to melting, and calving. Melting can be due to solar radiation, or heat and salt transport from the ocean. Iceberg calving is generally enhanced by waves impacting the iceberg.
Melting tends to be driven by the ocean, rather than solar radiation. Ocean driven melting is often modelled as
where is the melt rate in m/day, is the relative velocity between the iceberg and the ocean, is the temperature difference between the ocean and the iceberg, and is the length of the iceberg. is a constant based on properties of the iceberg and the ocean and is approximately in the polar ocean.
The influence of the shape of an iceberg and of the Coriolis force on iceberg melting rates has been demonstrated in laboratory experiments.
Wave erosion is more poorly constrained but can be estimated by
where is the wave erosion rate in m/day, , describes the sea state, is the sea surface temperature, and is the sea ice concentration.
Bubbles
Air trapped in snow forms bubbles as the snow is compressed to form firn and then glacial ice. Icebergs can contain up to 10% air bubbles by volume. These bubbles are released during melting, producing a fizzing sound that some may call "Bergie Seltzer". This sound results when the water-ice interface reaches compressed air bubbles trapped in the ice. As each bubble bursts it makes a "popping" sound and the acoustic properties of these bubbles can be used to study iceberg melt.
Stability
An iceberg may flip, or capsize, as it melts and breaks apart, changing the center of gravity. Capsizing can occur shortly after calving when the iceberg is young and establishing balance. Icebergs are unpredictable and can capsize anytime and without warning. Large icebergs that break off from a glacier front and flip onto the glacier face can push the entire glacier backwards momentarily, producing 'glacial earthquakes' that generate as much energy as an atomic bomb.
Color
Icebergs are generally white because they are covered in snow, but can be green, blue, yellow, black, striped, or even rainbow-colored. Seawater, algae and lack of air bubbles in the ice can create diverse colors. Sediment can create the dirty black coloration present in some icebergs.
Shape
In addition to size classification (Table 1), icebergs can be classified on the basis of their shapes. The two basic types of iceberg forms are tabular and non-tabular. Tabular icebergs have steep sides and a flat top, much like a plateau, with a length-to-height ratio of more than 5:1.
This type of iceberg, also known as an ice island, can be quite large, as in the case of Pobeda Ice Island. Antarctic icebergs formed by breaking off from an ice shelf, such as the Ross Ice Shelf or Filchner–Ronne Ice Shelf, are typically tabular. The largest icebergs in the world are formed this way.
Non-tabular icebergs have different shapes and include:
Dome: An iceberg with a rounded top.
Pinnacle: An iceberg with one or more spires.
Wedge: An iceberg with a steep edge on one side and a slope on the opposite side.
Dry-dock: An iceberg that has eroded to form a slot or channel.
Blocky: An iceberg with steep, vertical sides and a flat top. It differs from tabular icebergs in that its aspect ratio, the ratio between its width and height, is small, more like that of a block than a flat sheet.
Monitoring and control
History
Prior to 1914 there was no system in place to track icebergs to guard ships against collisions despite fatal sinkings of ships by icebergs. In 1907, SS Kronprinz Wilhelm, a German liner, rammed an iceberg and suffered a crushed bow, but she was still able to complete her voyage. The advent of watertight compartmentalization in ship construction led designers to declare their ships "unsinkable".
During the 1912 sinking of the Titanic, the iceberg that sank the Titanic killed more than 1,500 of its estimated 2,224 passengers and crew, seriously damaging the 'unsinkable' claim. For the remainder of the ice season of that year, the United States Navy patrolled the waters and monitored ice movements. In November 1913, the International Conference on the Safety of Life at Sea met in London to devise a more permanent system of observing icebergs. Within three months the participating maritime nations had formed the International Ice Patrol (IIP). The goal of the IIP was to collect data on meteorology and oceanography to measure currents, ice-flow, ocean temperature, and salinity levels. They monitored iceberg dangers near the Grand Banks of Newfoundland and provided the "limits of all known ice" in that vicinity to the maritime community. The IIP published their first records in 1921, which allowed for a year-by-year comparison of iceberg movement.
Technological development
Aerial surveillance of the seas in the early 1930s allowed for the development of charter systems that could accurately detail the ocean currents and iceberg locations. In 1945, experiments tested the effectiveness of radar in detecting icebergs. A decade later, oceanographic monitoring outposts were established for the purpose of collecting data; these outposts continue to serve in environmental study. A computer was first installed on a ship for the purpose of oceanographic monitoring in 1964, which allowed for a faster evaluation of data. By the 1970s, ice-breaking ships were equipped with automatic transmissions of satellite photographs of ice in Antarctica. Systems for optical satellites had been developed but were still limited by weather conditions. In the 1980s, drifting buoys were used in Antarctic waters for oceanographic and climate research. They are equipped with sensors that measure ocean temperature and currents.
Side looking airborne radar (SLAR) made it possible to acquire images regardless of weather conditions. On November 4, 1995, Canada launched RADARSAT-1. Developed by the Canadian Space Agency, it provides images of Earth for scientific and commercial purposes. This system was the first to use synthetic aperture radar (SAR), which sends microwave energy to the ocean surface and records the reflections to track icebergs. The European Space Agency launched ENVISAT (an observation satellite that orbits the Earth's poles) on March 1, 2002. ENVISAT employs advanced synthetic aperture radar (ASAR) technology, which can detect changes in surface height accurately. The Canadian Space Agency launched RADARSAT-2 in December 2007, which uses SAR and multi-polarization modes and follows the same orbit path as RADARSAT-1.
Modern monitoring
Iceberg concentrations and size distributions are monitored worldwide by the U.S. National Ice Center (NIC), established in 1995, which produces analyses and forecasts of Arctic, Antarctic, Great Lakes and Chesapeake Bay ice conditions. More than 95% of the data used in its sea ice analyses are derived from the remote sensors on polar-orbiting satellites that survey these remote regions of the Earth.
The NIC is the only organization that names and tracks all Antarctic Icebergs. It assigns each iceberg larger than along at least one axis a name composed of a letter indicating its point of origin and a running number. The letters used are as follows:
A – longitude 0° to 90° W (Bellingshausen Sea, Weddell Sea)
B – longitude 90° W to 180° (Amundsen Sea, Eastern Ross Sea)
C – longitude 90° E to 180° (Western Ross Sea, Wilkes Land)
D – longitude 0° to 90° E (Amery Ice Shelf, Eastern Weddell Sea)
The Danish Meteorological Institute monitors iceberg populations around Greenland using data collected by the synthetic aperture radar (SAR) on the Sentinel-1 satellites.
Iceberg management
In Labrador and Newfoundland, iceberg management plans have been developed to protect offshore installations from impacts with icebergs.
Commercial use
The idea of towing large icebergs to other regions as a source of water has been raised since at least the 1950s, without having been put into practice. In 2017, a business from the UAE announced plans to tow an iceberg from Antarctica to the Middle East; in 2019 salvage engineer Nick Sloane announced a plan to move one to South Africa at an estimated cost of $200 million. In 2019, a German company, Polewater, announced plans to tow Antarctic icebergs to places like South Africa.
Companies have used iceberg water in products such as bottled water, fizzy ice cubes and alcoholic drinks. For example, Iceberg Beer by Quidi Vidi Brewing Company is made from icebergs found around St. John's, Newfoundland. Although annual iceberg supply in Newfoundland and Labrador exceeds the total freshwater consumption of the United States, in 2016 the province introduced a tax on iceberg harvesting and imposed a limit on how much fresh water can be exported yearly.
Oceanography and ecology
The freshwater injected into the ocean by melting icebergs can change the density of the seawater in the vicinity of the iceberg. Fresh melt water released at depth is lighter, and therefore more buoyant, than the surrounding seawater causing it to rise towards the surface. Icebergs can also act as floating breakwaters, impacting ocean waves.
Icebergs contain variable concentrations of nutrients and minerals that are released into the ocean during melting. Iceberg-derived nutrients, particularly the iron contained in sediments, can fuel blooms of phytoplankton. Samples collected from icebergs in Antarctica, Patagonia, Greenland, Svalbard, and Iceland, however, show that iron concentrations vary significantly, complicating efforts to generalize the impacts of icebergs on marine ecosystems.
Recent large icebergs
Iceberg B15 calved from the Ross Ice Shelf in 2000 and initially had an area of . It broke apart in November 2002. The largest remaining piece of it, Iceberg B-15A, with an area of , was still the largest iceberg on Earth until it ran aground and split into several pieces October 27, 2005, an event that was observed by seismographs both on the iceberg and across Antarctica. It has been hypothesized that this breakup may also have been abetted by ocean swell generated by an Alaskan storm 6 days earlier and away.
1987, Iceberg B-9,
1998, Iceberg A-38, about
1999, Iceberg B-17B , shipping alert issued December 2009.
2000, Iceberg B-15
2002, Iceberg C-19,
2002, Iceberg B-22,
2003 broke off, Iceberg B-15A,
2006, Iceberg D-16,
2010, Ice sheet, , broken off of Petermann Glacier in northern Greenland on August 5, 2010, considered to be the largest Arctic iceberg since 1962. About a month later, this iceberg split into two pieces upon crashing into Joe Island in the Nares Strait next to Greenland. In June 2011, large fragments of the Petermann Ice Islands were observed off the Labrador coast.
2014, Iceberg B-31, , 2014
2017, Iceberg A-68, (Larsen C)
2018, Iceberg B-46,
2019, Iceberg D-28,
2021, Iceberg A-74 from the Brunt Ice Shelf,
2021, Iceberg A-76 from the Ronne Ice Shelf,
In culture
One of the most infamous icebergs in history is the iceberg that sank the Titanic. The catastrophe led to the establishment of an International Ice Patrol shortly after. Icebergs in both the northern and southern hemispheres have often been compared in size to multiples of the -area of Manhattan Island.
Artists have used icebergs as the subject matter for their paintings. Frederic Edwin Church, The Icebergs, 1861 was painted from sketches Church completed on a boat trip off Newfoundland and Labrador. Caspar David Friedrich, The Sea of Ice, 1823–1824 is a polar landscape with an iceberg and ship wreck depicting the dangers of such conditions. William Bradford created detailed paintings of sailing ships set in arctic coasts and was fascinated by icebergs. Albert Bierstadt made studies on arctic trips aboard steamships in 1883 and 1884 that were the basis of his paintings of arctic scenes with colossal icebergs made in the studio.
American poet, Lydia Sigourney, wrote the poem "Icebergs". While on a return journey from Europe in 1841, her steamship encountered a field of icebergs overnight, during an Aurora Borealis. The ship made it through unscathed to the next morning, when the sun rose and "touched the crowns, Of all those arctic kings."
| Physical sciences | Glaciology | null |
146237 | https://en.wikipedia.org/wiki/Osprey | Osprey | The osprey (; Pandion haliaetus), historically known as sea hawk, river hawk, and fish hawk, is a diurnal, fish-eating bird of prey with a cosmopolitan range. It is a large raptor, reaching more than in length and across the wings. It is brown on the upperparts and predominantly greyish on the head and underparts.
The osprey tolerates a wide variety of habitats, nesting in any location near a body of water providing an adequate food supply. It is found on all continents except Antarctica, although in South America it occurs only as a non-breeding migrant.
As its other common names suggest, the osprey's diet consists almost exclusively of fish. It possesses specialised physical characteristics and unique behaviour in hunting its prey. Because of its unique characteristics it is classified in its own taxonomic genus, Pandion, and family, Pandionidae.
Taxonomy
The osprey was described in 1758 by the Swedish naturalist Carl Linnaeus under the name Falco haliaetus in his landmark tenth edition of his Systema Naturae. Linnaeus specified the type locality as Europe, but in 1761 he restricted the locality to Sweden. The osprey is the only species placed in the genus Pandion that was introduced by the French zoologist Marie Jules César Savigny in 1809. The genus is the sole member of the family Pandionidae. The species has always presented a riddle to taxonomists, but here it is treated as the sole living member of the family Pandionidae, and the family listed in its traditional place as part of the order Accipitriformes. Other schemes place it alongside the hawks and eagles in the family Accipitridae. The Sibley-Ahlquist taxonomy has placed it together with the other diurnal raptors in a greatly enlarged Ciconiiformes, but this results in an unnatural paraphyletic classification. Molecular phylogenetic analysis has found that the family Pandionidae is sister to the family Accipitridae. It is estimated that the two families diverged around 50.8 million years ago.The osprey is unusual in that it is a sole living species that occurs nearly worldwide. Even the few subspecies are not unequivocally separable. There are four generally recognised subspecies, although differences are small, and ITIS lists only the first three.
Pandion haliaetus haliaetus (Linnaeus, 1758) – the Eurasian osprey is the nominate subspecies that occurs across the Palearctic realm and several parts of sub-Saharan Africa from the Azores and the Iberian Peninsula east to Japan and Kamchatka Peninsula, throughout South and Southeast Asia, the Indian subcontinent, Madagascar and much of the African coastline.
P. haliaetus carolinensis (Gmelin, 1788)– the American or North American osprey occurs from Alaska and Canada to much of Central and South America, except Chile and Patagonia. It is larger, darker bodied and has a paler breast than the European osprey.
P. haliaetus ridgwayi Maynard, 1887 – Ridgway's osprey occurs in the Caribbean islands. It has a very pale head and breast and a weak eye mask. It is non-migratory. Its scientific name commemorates Robert Ridgway.
P. haliaetus cristatus (Vieillot, 1816) – the Australasian osprey is the smallest and most distinctive subspecies that occurs along the entire marine coastline of Australia and some larger freshwater rivers as well as in Tasmania. It is not migratory. Some authorities have assigned it full species-status as Pandion cristatus, also known as the eastern osprey. A 2018 genetic study using microsatellite data showed only low genetic divergence between cristatus and the other subspecies.
Fossil record
Two extinct species were named from the fossil record. Pandion homalopteron described by Stuart L. Warter in 1976 was found in marine Middle Miocene deposits of the Barstovian age in the southern part of California. The second species Pandion lovensis was described by Jonathan J. Becker in 1985 and found in Florida; it dates to the Late Clarendonian and possibly represents a separate lineage from that of P. homalopteron and P. haliaetus. A number of claw fossils have been recovered from Pliocene and Pleistocene sediments in Florida and South Carolina.
The oldest recognized family Pandionidae fossils were recovered from the Oligocene age Jebel Qatrani Formation in Faiyum Governorate, Egypt. However, they are not complete enough to assign to a specific genus. Another Pandionidae claw fossil was recovered from Early Oligocene deposits in the Mainz basin, Germany, and was described in 2006 by Gerald Mayr.
Etymology
The genus name Pandion derives from Pandíōn , the mythical Greek king of Athens and grandfather of Theseus, Pandion II. The species name haliaetus () comes from Greek haliáetos "sea-eagle" (also haliaietos) from the combining form hali- of hals "sea" and aetos, "eagle".
The origins of osprey are obscure; the word itself was first recorded around 1460, derived via the Anglo-French ospriet and the Medieval Latin avis prede "bird of prey," from the Latin avis praedae though the Oxford English Dictionary notes a connection with the Latin ossifraga or "bone breaker" of Pliny the Elder. However, this term referred to the bearded vulture.
Description
The osprey differs in several respects from other diurnal birds of prey. Its toes are of equal length, its tarsi are reticulate, and its talons are rounded, rather than grooved. The osprey and owls are the only raptors whose outer toe is reversible, allowing them to grasp their prey with two toes in front and two behind. This is particularly helpful when they grab slippery fish.
The osprey is in weight and in length with a wingspan. It is, thus, of similar size to the largest members of the Buteo or Falco genera. The subspecies are fairly close in size, with the nominate subspecies averaging , P. h. carolinensis averaging and P. h. cristatus averaging . The wing chord measures , the tail measures and the tarsus is .
The upperparts are a deep, glossy brown, while the breast is white, sometimes streaked with brown, and the underparts are pure white. The head is white with a dark mask across the eyes, reaching to the sides of the neck. The irises of the eyes are golden to brown, and the transparent nictitating membrane is pale blue. The bill is black, with a blue cere, and the feet are white with black talons. On the underside of the wings the wrists are black, which serves as a field mark. A short tail and long, narrow wings with four long, finger-like feathers, and a shorter fifth, give it a very distinctive appearance.
The sexes appear fairly similar, but the adult male can be distinguished from the female by its slimmer body and narrower wings. The breast band of the male is also weaker than that of the female or is non-existent, and the underwing coverts of the male are more uniformly pale. It is straightforward to determine the sex in a breeding pair, but harder with individual birds.
The juvenile osprey may be identified by buff fringes to the plumage of the upperparts, a buff tone to the underparts, and streaked feathers on the head. During spring, barring on the underwings and flight feathers is a better indicator of a young bird, due to wear on the upperparts.
In flight, the osprey has arched wings and drooping "hands", giving it a gull-like appearance. The call is a series of sharp whistles, described as cheep, cheep, or yewk, yewk. If disturbed by activity near the nest, the call is a frenzied cheereek!
Distribution and habitat
The osprey is the second most widely distributed raptor species, after the peregrine falcon, and is one of only six land-birds with a worldwide distribution. It is found in temperate and tropical regions of all continents, except Antarctica. In North America it breeds from Alaska and Newfoundland south to the Gulf Coast and Florida, wintering further south from the southern United States through to Argentina. It is found in summer throughout Europe north into Ireland, Scandinavia, Finland and Great Britain though not Iceland, and winters in North Africa. In Australia it is mainly sedentary and found patchily around the coastline, though it is a non-breeding visitor to eastern Victoria and Tasmania.
There is a gap, corresponding with the coast of the Nullarbor Plain, between its westernmost breeding site in South Australia and the nearest breeding sites to the west in Western Australia. In the islands of the Pacific it is found in the Bismarck Islands, Solomon Islands and New Caledonia, and fossil remains of adults and juveniles have been found in Tonga, where it probably was wiped out by arriving humans. It is possible it may once have ranged across Vanuatu and Fiji as well. It is an uncommon to fairly common winter visitor to all parts of South Asia, and Southeast Asia from Myanmar through to Indochina and southern China, Indonesia, Malaysia, and the Philippines.
Behaviour and ecology
Diet
The osprey is piscivorous, with fish making up 99% of its diet. It typically takes live fish weighing and about in length, but virtually any type of fish from to can be taken. Even larger northern pike (Esox lucius) has been taken in Russia. The species rarely scavenges dead or dying fish.
Ospreys have a vision that is well adapted to detecting underwater objects from the air. Prey is first sighted when the osprey is above the water, after which the bird hovers momentarily and then plunges feet first into the water. They catch fish by diving into a body of water, oftentimes completely submerging their entire bodies. As an osprey dives it adjusts the angle of its flight to account for the distortion of the fish's image caused by refraction. Ospreys will typically eat on a nearby perch but have also been known to carry fish for longer distances.
Occasionally, the osprey may prey on rodents, rabbits, hares, other mammals, snakes, turtles, frogs, birds, salamanders, conchs, and crustaceans. Reports of ospreys feeding on carrion are rare. They have been observed eating dead white-tailed deer and Virginia opossums.
Adaptations
The osprey has several adaptations that suit its piscivorous lifestyle. These include reversible outer toes, sharp spicules on the underside of the toes, closable nostrils to keep out water during dives, backward-facing scales on the talons which act as barbs to help hold its catch and dense plumage which is oily and prevents its feathers from getting waterlogged.
Reproduction
The osprey breeds near freshwater lakes and rivers, and sometimes on coastal brackish waters. Rocky outcrops just offshore are used in Rottnest Island off the coast of Western Australia, where there are 14 or so similar nesting sites of which five to seven are used in any one year. Many are renovated each season, and some have been used for 70 years. The nest is a large heap of sticks, driftwood, turf, or seaweed built in forks of trees, rocky outcrops, utility poles, artificial platforms, or offshore islets. As wide as 2 meters and weighing about , large nests on utility poles may be fire hazards and have caused power outages.
Generally, ospreys reach sexual maturity and begin breeding around the age of three to four, though in some regions with high osprey densities, such as Chesapeake Bay in the United States, they may not start breeding until five to seven years old, and there may be a shortage of suitable tall structures. If there are no nesting sites available, young ospreys may be forced to delay breeding. To ease this problem, posts are sometimes erected to provide more sites suitable for nest building. The nesting platform design developed by the organization Citizens United to Protect the Maurice River and Its Tributaries, Inc. has become the official design of the State of New Jersey, U.S. The nesting platform plans and materials list, available online, have been utilized by people from a number of different geographical regions. There is a global site for mapping osprey nest locations and logging observations on reproductive success.
Ospreys usually mate for life. Rarely, polyandry has been recorded. The breeding season varies according to latitude: spring (September–October) in southern Australia, April to July in northern Australia, and winter (June–August) in southern Queensland. In spring, the pair begins a five-month period of partnership to raise their young. The female lays two to four eggs within a month and relies on the size of the nest to conserve heat. The eggs are whitish with bold splotches of reddish-brown and are about and weigh about . The eggs are incubated for about 35–43 days to hatching.
The newly hatched chicks weigh only , but fledge in 8–10 weeks. A study on Kangaroo Island, South Australia, had an average time between hatching and fledging of 69 days. The same study found an average of 0.66 young fledged per year per occupied territory, and 0.92 young fledged per year per active nest. Some 22% of surviving young either remained on the island or returned at maturity to join the breeding population. When food is scarce, the first chicks to hatch are most likely to survive. The typical lifespan is 7–10 years, though rarely individuals can grow to as old as 20–25 years. The oldest European wild osprey on record lived to be over thirty years of age.
Migration
European breeders winter in Africa. American and Canadian breeders winter in South America, although some stay in the southernmost U.S. states such as Florida and California. Some ospreys from Florida migrate to South America. Australasian ospreys tend not to migrate.
Studies of Swedish ospreys showed that females tend to migrate to Africa earlier than males. More stopovers are made during their autumn migration. The variation of timing and duration in autumn was more variable than in spring. Although migrating predominantly during the day, they sometimes fly in the dark hours, particularly in crossings over water and cover on average per day with a maximum of per day. European birds may also winter in South Asia, as indicated by an osprey tagged in Norway being monitored in western India. In the Mediterranean, ospreys show partial migratory behaviour with some individuals remaining resident, whilst others undertake relatively short migration trips.
Mortality
Swedish ospreys have a significantly higher mortality rate during migration seasons than during stationary periods, with more than half of the total annual mortality occurring during migration. These deaths can also be categorized into spatial patterns: Spring mortality occurs mainly in Africa, which can be traced to crossing the Sahara desert. Mortality can also occur through mishaps with human utilities, such as nesting near overhead electric cables or collisions with aircraft.
Conservation
The osprey has a large range, covering in just Africa and the Americas, and has a large global population estimated at 460,000 individuals. Although global population trends have not been quantified, the species is not believed to approach the thresholds for the population decline criterion of the IUCN Red List (i.e., declining more than 30% in ten years or three generations), and for these reasons, the species is evaluated as Least Concern. There is evidence for regional decline in South Australia where former territories at locations in the Spencer Gulf and along the lower Murray River have been vacant for decades.
In the late 19th and early 20th centuries, the main threats to osprey populations were egg collectors and hunting of the adults along with other birds of prey, but osprey populations declined drastically in many areas in the 1950s and 1960s; this appeared to be in part due to the toxic effects of insecticides such as DDT on reproduction. The pesticide interfered with the bird's calcium metabolism which resulted in thin-shelled, easily broken or infertile eggs. Possibly because of the banning of DDT in many countries in the early 1970s, together with reduced persecution, the osprey, as well as other affected bird of prey species, have made significant recoveries. In South Australia, nesting sites on the Eyre Peninsula and Kangaroo Island are vulnerable to unmanaged coastal recreation and encroaching urban development.
Cultural depictions
Literature
The Roman writer Pliny the Elder reported that parent ospreys made their young fly up to the sun as a test, and dispatched any that failed.
Another odd legend regarding this fish-eating bird of prey, derived from the writings of Albertus Magnus and recorded in Holinshed's Chronicles, was that it had one webbed foot and one taloned foot.
The osprey is mentioned in the famous Chinese folk poem "guan guan ju jiu" (關關雎鳩); "ju jiu" 雎鳩 refers to the osprey, and "guan guan" (關關) to its voice. In the poem, the osprey is considered to be an icon of fidelity and harmony between wife and husband, due to its highly monogamous habits. Some commentators have claimed that "ju jiu" in the poem is not the osprey but the mallard duck, since the osprey cannot make the sound "guan guan".
The Irish poet William Butler Yeats used a grey wandering osprey as a representation of sorrow in The Wanderings of Oisin and Other Poems (1889).
There was a medieval belief that fish were so mesmerised by the osprey that they turned belly-up in surrender, and this is referenced by Shakespeare in Act 4 Scene 5 of Coriolanus:
Iconography
In heraldry, the osprey is typically depicted as a white eagle, often maintaining a fish in its talons or beak, and termed a "sea-eagle". It is historically regarded as a symbol of vision and abundance; more recently it has become a symbol of positive responses to nature, and has been featured on more than 50 international postage stamps.
In 1994, the osprey was declared the provincial bird of Nova Scotia, Canada.
Sports
Some sports clubs are named after the osprey such as the University of North Florida's North Florida Ospreys and Missoula Osprey baseball team.
"Seahawks", another term for osprey, is also common among sports teams. The Seattle Seahawks, a professional American football team in the National Football League, received their identity from a naming contest, defeating 1,740 others. According to team general manager John Thompson, the name "shows aggressiveness, reflects our soaring Northwest heritage and belongs to no other major league team."
Other
So-called "osprey" plumes were an important item in the plume trade of the late 19th century and used in hats including those used as part of the army uniform. Despite their name, these plumes were actually obtained from egrets.
During the 2017 regular session of the Oregon Legislature, there was a short-lived controversy over the western meadowlark's status as the state bird versus the osprey. The sometimes-spirited debate included state representative Rich Vial playing the meadowlark's song on his smartphone over the House microphone. A compromise was reached in SCR 18, which was passed on the last day of the session, designating the western meadowlark as the state songbird and the osprey as the state raptor.
| Biology and health sciences | Accipitriformes and Falconiformes | null |
146249 | https://en.wikipedia.org/wiki/Earth%27s%20outer%20core | Earth's outer core | Earth's outer core is a fluid layer about thick, composed of mostly iron and nickel that lies above Earth's solid inner core and below its mantle. The outer core begins approximately beneath Earth's surface is at the core-mantle boundary and ends beneath Earth's surface at the inner core boundary.
Properties
The outer core of Earth is liquid, unlike its inner core, which is solid. Evidence for a fluid outer core includes seismology which shows that seismic shear-waves are not transmitted through the outer core. Although having a composition similar to Earth's solid inner core, the outer core remains liquid as there is not enough pressure to keep it in a solid state.
Seismic inversions of body waves and normal modes constrain the radius of the outer core to be 3483 km with an uncertainty of 5 km, while that of the inner core is 1220±10 km.
Estimates for the temperature of the outer core are about in its outer region and near the inner core. Modeling has shown that the outer core, because of its high temperature, is a low-viscosity fluid that convects turbulently. The dynamo theory sees eddy currents in the nickel-iron fluid of the outer core as the principal source of Earth's magnetic field. The average magnetic field strength in Earth's outer core is estimated to be 2.5 millitesla, 50 times stronger than the magnetic field at the surface.
As Earth's core cools, the liquid at the inner core boundary freezes, causing the solid inner core to grow at the expense of the outer core, at an estimated rate of 1 mm per year. This is approximately 80,000 tonnes of iron per second.
Light elements of Earth's outer core
Composition
Earth's outer core cannot be entirely constituted of iron or iron-nickel alloy because their densities are higher than geophysical measurements of the density of Earth's outer core. In fact, Earth's outer core is approximately 5 to 10 percent lower density than iron at Earth's core temperatures and pressures. Hence it has been proposed that light elements with low atomic numbers compose part of Earth's outer core, as the only feasible way to lower its density. Although Earth's outer core is inaccessible to direct sampling, the composition of light elements can be meaningfully constrained by high-pressure experiments, calculations based on seismic measurements, models of Earth's accretion, and carbonaceous chondrite meteorite comparisons with bulk silicate Earth (BSE). Recent estimates are that Earth's outer core is composed of iron along with 0 to 0.26 percent hydrogen, 0.2 percent carbon, 0.8 to 5.3 percent oxygen, 0 to 4.0 percent silicon, 1.7 percent sulfur, and 5 percent nickel by weight, and the temperature of the core-mantle boundary and the inner core boundary ranges from 4,137 to 4,300 K and from 5,400 to 6,300 K respectively.
Constraints
Accretion
The variety of light elements present in Earth's outer core is constrained in part by Earth's accretion. Namely, the light elements contained must have been abundant during Earth's formation, must be able to partition into liquid iron at low pressures, and must not volatilize and escape during Earth's accretionary process.
CI chondrites
CI chondritic meteorites are believed to contain the same planet-forming elements in the same proportions as in the early Solar System, so differences between CI meteorites and BSE can provide insights into the light element composition of Earth's outer core. For instance, the depletion of silicon in BSE compared to CI meteorites may indicate that silicon was absorbed into Earth's core; however, a wide range of silicon concentrations in Earth's outer and inner core is still possible.
Implications for Earth's accretion and core formation history
Tighter constraints on the concentrations of light elements in Earth's outer core would provide a better understanding of Earth's accretion and core formation history.
Consequences for Earth's accretion
Models of Earth's accretion could be better tested if we had better constraints on light element concentrations in Earth's outer core. For example, accretionary models based on core-mantle element partitioning tend to support proto-Earths constructed from reduced, condensed, and volatile-free material, despite the possibility that oxidized material from the outer Solar System was accreted towards the conclusion of Earth's accretion. If we could better constrain the concentrations of hydrogen, oxygen, and silicon in Earth's outer core, models of Earth's accretion that match these concentrations would presumably better constrain Earth’s formation.
Consequences for Earth's core formation
The depletion of siderophile elements in Earth's mantle compared to chondritic meteorites is attributed to metal-silicate reactions during formation of Earth's core. These reactions are dependent on oxygen, silicon, and sulfur, so better constraints on concentrations of these elements in Earth's outer core will help elucidate the conditions of formation of Earth's core.
In another example, the possible presence of hydrogen in Earth's outer core suggests that the accretion of Earth’s water was not limited to the final stages of Earth's accretion and that water may have been absorbed into core-forming metals through a hydrous magma ocean.
Implications for Earth's magnetic field
Earth's magnetic field is driven by thermal convection and also by chemical convection, the exclusion of light elements from the inner core, which float upward within the fluid outer core while denser elements sink. This chemical convection releases gravitational energy that is then available to power the geodynamo that produces Earth's magnetic field. Carnot efficiencies with large uncertainties suggest that compositional and thermal convection contribute about 80 percent and 20 percent respectively to the power of Earth's geodynamo. Traditionally it was thought that prior to the formation of Earth's inner core, Earth's geodynamo was mainly driven by thermal convection. However, recent claims that the thermal conductivity of iron at core temperatures and pressures is much higher than previously thought imply that core cooling was largely by conduction not convection, limiting the ability of thermal convection to drive the geodynamo. This conundrum is known as the new "core paradox." An alternative process that could have sustained Earth's geodynamo requires Earth's core to have initially been hot enough to dissolve oxygen, magnesium, silicon, and other light elements. As the Earth's core began to cool, it would become supersaturated in these light elements that would then precipitate into the lower mantle forming oxides leading to a different variant of chemical convection.
The magnetic field generated by core flow is essential to protect life from interplanetary radiation and prevent the atmosphere from dissipating in the solar wind. The rate of cooling by conduction and convection is uncertain, but one estimate is that the core would not be expected to freeze up for approximately 91 billion years, which is well after the Sun is expected to expand, sterilize the surface of the planet, and then burn out.
| Physical sciences | Geophysics | null |
146250 | https://en.wikipedia.org/wiki/Carbon-14 | Carbon-14 | Carbon-14, C-14, C or radiocarbon, is a radioactive isotope of carbon with an atomic nucleus containing 6 protons and 8 neutrons. Its presence in organic matter is the basis of the radiocarbon dating method pioneered by Willard Libby and colleagues (1949) to date archaeological, geological and hydrogeological samples. Carbon-14 was discovered on February 27, 1940, by Martin Kamen and Sam Ruben at the University of California Radiation Laboratory in Berkeley, California. Its existence had been suggested by Franz Kurie in 1934.
There are three naturally occurring isotopes of carbon on Earth: carbon-12 (C), which makes up 99% of all carbon on Earth; carbon-13 (C), which makes up 1%; and carbon-14 (C), which occurs in trace amounts, making up about 1-1.5 atoms per 10 atoms of carbon in the atmosphere. C and C are both stable; C is unstable, with half-life years. Carbon-14 has a specific activity of 62.4 mCi/mmol (2.31 GBq/mmol), or 164.9 GBq/g. Carbon-14 decays into nitrogen-14 () through beta decay. A gram of carbon containing 1 atom of carbon-14 per 10 atoms, emits ~0.2 beta (β) particles per second. The primary natural source of carbon-14 on Earth is cosmic ray action on nitrogen in the atmosphere, and it is therefore a cosmogenic nuclide. However, open-air nuclear testing between 1955 and 1980 contributed to this pool.
The different isotopes of carbon do not differ appreciably in their chemical properties. This resemblance is used in chemical and biological research, in a technique called carbon labeling: carbon-14 atoms can be used to replace nonradioactive carbon, in order to trace chemical and biochemical reactions involving carbon atoms from any given organic compound.
Radioactive decay and detection
Carbon-14 undergoes beta decay:
→ + + + 156.5 keV
By emitting an electron and an electron antineutrino, one of the neutrons in carbon-14 decays to a proton and the carbon-14 (half-life of years) decays into the stable (non-radioactive) isotope nitrogen-14.
As usual with beta decay, almost all the decay energy is carried away by the beta particle and the neutrino. The emitted beta particles have a maximum energy of about 156 keV, while their weighted mean energy is 49 keV. These are relatively low energies; the maximum distance traveled is estimated to be 22 cm in air and 0.27 mm in body tissue. The fraction of the radiation transmitted through the dead skin layer is estimated to be 0.11. Small amounts of carbon-14 are not easily detected by typical Geiger–Müller (G-M) detectors; it is estimated that G-M detectors will not normally detect contamination of less than about 100,000 decays per minute (0.05 μCi). Liquid scintillation counting is the preferred method although more recently, accelerator mass spectrometry has become the method of choice; it counts all the carbon-14 atoms in the sample and not just the few that happen to decay during the measurements; it can therefore be used with much smaller samples (as small as individual plant seeds), and gives results much more quickly. The G-M counting efficiency is estimated to be 3%. The half-distance layer in water is 0.05 mm.
Radiocarbon dating
Radiocarbon dating is a radiometric dating method that uses C to determine the age of carbonaceous materials up to about 60,000 years old. The technique was developed by Willard Libby and his colleagues in 1949 during his tenure as a professor at the University of Chicago. Libby estimated that the radioactivity of exchangeable C would be about 14 decays per minute (dpm) per gram of carbon, and this is still used as the activity of the modern radiocarbon standard. In 1960, Libby was awarded the Nobel Prize in chemistry for this work.
One of the frequent uses of the technique is to date organic remains from archaeological sites. Plants fix atmospheric carbon during photosynthesis; so the level of C in plants and animals when they die, roughly equals the level of C in the atmosphere at that time. However, it thereafter decreases exponentially; so the date of death or fixation can be estimated. The initial C level for the calculation can either be estimated, or else directly compared with known year-by-year data from tree-ring data (dendrochronology) up to 10,000 years ago (using overlapping data from live and dead trees in a given area), or else from cave deposits (speleothems), back to about 45,000 years before present. A calculation or (more accurately) a direct comparison of carbon-14 levels in a sample, with tree ring or cave-deposit C levels of a known age, then gives the wood or animal sample age-since-formation. Radiocarbon is also used to detect disturbance in natural ecosystems; for example, in peatland landscapes, radiocarbon can indicate that carbon which was previously stored in organic soils is being released due to land clearance or climate change.
Cosmogenic nuclides are also used as proxy data to characterize cosmic particle and solar activity of the distant past.
Origin
Natural production in the atmosphere
Carbon-14 is produced in the upper troposphere and the stratosphere by thermal neutrons absorbed by nitrogen atoms. When cosmic rays enter the atmosphere, they undergo various transformations, including the production of neutrons. The resulting neutrons (n) participate in the following n-p reaction (p is proton):
+ n → + p + 0.626 MeV
The highest rate of carbon-14 production takes place at altitudes of and at high geomagnetic latitudes.
The rate of C production can be modeled, yielding values of 16,400 or 18,800 atoms of per second per square meter of Earth's surface, which agrees with the global carbon budget that can be used to backtrack, but attempts to measure the production time directly in situ were not very successful. Production rates vary because of changes to the cosmic ray flux caused by the heliospheric modulation (solar wind and solar magnetic field), and, of great significance, due to variations in the Earth's magnetic field. Changes in the carbon cycle however can make such effects difficult to isolate and quantify.
Occasional spikes may occur; for example, there is evidence for an unusually high production rate in AD 774–775, caused by an extreme
solar energetic particle event, the strongest such event to have occurred within the last ten millennia. Another "extraordinarily large" C increase (2%) has been associated with a 5480 BC event, which is unlikely to be a solar energetic particle event.
Carbon-14 may also be produced by lightning but in amounts negligible, globally, compared to cosmic ray production. Local effects of cloud-ground discharge through sample residues are unclear, but possibly significant.
Other carbon-14 sources
Carbon-14 can also be produced by other neutron reactions, including in particular C(n,γ)C and O(n,α)C with thermal neutrons, and N(n,d)C and O(n,He)C with fast neutrons. The most notable routes for C production by thermal neutron irradiation of targets (e.g., in a nuclear reactor) are summarized in the table.
Another source of carbon-14 is cluster decay branches from traces of naturally occurring isotopes of radium, though this decay mode has a branching ratio on the order of relative to alpha decay, so radiogenic carbon-14 is extremely rare.
Formation during nuclear tests
The above-ground nuclear tests that occurred in several countries in 1955-1980 (see List of nuclear tests) dramatically increased the amount of C in the atmosphere and subsequently the biosphere; after the tests ended, the atmospheric concentration of the isotope began to decrease, as radioactive CO was fixed into plant and animal tissue, and dissolved in the oceans.
One side-effect of the change in atmospheric C is that this has enabled some options (e.g. bomb-pulse dating) for determining the birth year of an individual, in particular, the amount of C in tooth enamel, or the carbon-14 concentration in the lens of the eye.
In 2019, Scientific American reported that carbon-14 from nuclear testing has been found in animals from one of the most inaccessible regions on Earth, the Mariana Trench in the Pacific Ocean.
The concentration of C in atmospheric CO, reported as the C/C ratio with respect to a standard, has (since about 2022) declined to levels similar to those prior to the above-ground nuclear tests of the 1950s and 1960s. Though the extra C generated by those nuclear tests has not disappeared from the atmosphere, oceans and biosphere, it is diluted due to the Suess effect.
Emissions from nuclear power plants
Carbon-14 is produced in coolant at boiling water reactors (BWRs) and pressurized water reactors (PWRs). It is typically released into the air in the form of carbon dioxide at BWRs, and methane at PWRs. Best practice for nuclear power plant operator management of carbon-14 includes releasing it at night, when plants are not photosynthesizing. Carbon-14 is also generated inside nuclear fuels (some due to transmutation of oxygen in the uranium oxide, but most significantly from transmutation of nitrogen-14 impurities), and if the spent fuel is sent to nuclear reprocessing then the C is released, for example as CO during PUREX.
Occurrence
Dispersion in the environment
After production in the upper atmosphere, the carbon-14 reacts rapidly to form mostly (about 93%) CO (carbon monoxide), which subsequently oxidizes at a slower rate to form , radioactive carbon dioxide. The gas mixes rapidly and becomes evenly distributed throughout the atmosphere (the mixing timescale on the order of weeks). Carbon dioxide also dissolves in water and thus permeates the oceans, but at a slower rate. The atmospheric half-life for removal of has been estimated at roughly 12 to 16 years in the Northern Hemisphere. The transfer between the ocean shallow layer and the large reservoir of bicarbonates in the ocean depths occurs at a limited rate.
In 2009 the activity of was 238 Bq per kg carbon of fresh terrestrial biomatter, close to the values before atmospheric nuclear testing (226 Bq/kg C; 1950).
Total inventory
The inventory of carbon-14 in Earth's biosphere is about 300 megacuries (11 EBq), of which most is in the oceans.
The following inventory of carbon-14 has been given:
Global inventory: ~8500 PBq (about 50 t)
Atmosphere: 140 PBq (840 kg)
Terrestrial materials: the balance
From nuclear testing (until 1990): 220 PBq (1.3 t)
In fossil fuels
Many human-made chemicals are derived from fossil fuels (such as petroleum or coal) in which C is greatly depleted because the age of fossils far exceeds the half-life of C. The relative absence of is therefore used to determine the relative contribution (or mixing ratio) of fossil fuel oxidation to the total carbon dioxide in a given region of Earth's atmosphere.
Dating a specific sample of fossilized carbonaceous material is more complicated. Such deposits often contain trace amounts of C. These amounts can vary significantly between samples, ranging up to 1% of the ratio found in living organisms (an apparent age of about 40,000 years). This may indicate contamination by small amounts of bacteria, underground sources of radiation causing a N(n,p)C reaction, direct uranium decay (though reported measured ratios of C/U in uranium-bearing ores would imply roughly 1 uranium atom for every two carbon atoms in order to cause the C/C ratio, measured to be on the order of 10), or other unknown secondary sources of C production. The presence of C in the isotopic signature of a sample of carbonaceous material possibly indicates its contamination by biogenic sources or the decay of radioactive material in surrounding geologic strata. In connection with building the Borexino solar neutrino observatory, petroleum feedstock (for synthesizing the primary scintillant) was obtained with low C content. In the Borexino Counting Test Facility, a C/C ratio of 1.94×10 was determined; probable reactions responsible for varied levels of C in different petroleum reservoirs, and the lower C levels in methane, have been discussed by Bonvicini et al.
In the human body
Since many sources of human food are ultimately derived from terrestrial plants, the relative concentration of C in human bodies is nearly identical to the relative concentration in the atmosphere. The rates of disintegration of potassium-40 (K) and C in the normal adult body are comparable (a few thousand decays per second). The beta decays from external (environmental) radiocarbon contribute about 0.01 mSv/year (1 mrem/year) to each person's dose of ionizing radiation. This is small compared to the doses from K (0.39 mSv/year) and radon (variable).
C can be used as a radioactive tracer in medicine. In the initial variant of the urea breath test, a diagnostic test for Helicobacter pylori, urea labeled with about C is fed to a patient (i.e. 37,000 decays per second). In the event of a H. pylori infection, the bacterial urease enzyme breaks down the urea into ammonia and radioactively-labeled carbon dioxide, which can be detected by low-level counting of the patient's breath.
| Physical sciences | Group 14 | Chemistry |
146253 | https://en.wikipedia.org/wiki/Shock%20wave | Shock wave | In physics, a shock wave (also spelled shockwave), or shock, is a type of propagating disturbance that moves faster than the local speed of sound in the medium. Like an ordinary wave, a shock wave carries energy and can propagate through a medium, but is characterized by an abrupt, nearly discontinuous, change in pressure, temperature, and density of the medium.
For the purpose of comparison, in supersonic flows, additional increased expansion may be achieved through an expansion fan, also known as a Prandtl–Meyer expansion fan. The accompanying expansion wave may approach and eventually collide and recombine with the shock wave, creating a process of destructive interference. The sonic boom associated with the passage of a supersonic aircraft is a type of sound wave produced by constructive interference.
Unlike solitons (another kind of nonlinear wave), the energy and speed of a shock wave alone dissipates relatively quickly with distance. When a shock wave passes through matter, energy is preserved but entropy increases. This change in the matter's properties manifests itself as a decrease in the energy which can be extracted as work, and as a drag force on supersonic objects; shock waves are strongly irreversible processes.
Terminology
Shock waves can be:
Normal At 90° (perpendicular) to the shock medium's flow direction.
Oblique At an angle to the direction of flow.
Bow Occurs upstream of the front (bow) of a blunt object when the upstream flow velocity exceeds Mach 1.
Some other terms:
Shock front: The boundary over which the physical conditions undergo an abrupt change because of a shock wave.
Contact front: In a shock wave caused by a driver gas (for example the "impact" of a high explosive on the surrounding air), the boundary between the driver (explosive products) and the driven (air) gases. The contact front trails the shock front.
In supersonic flows
The abruptness of change in the features of the medium, that characterize shock waves, can be viewed as a phase transition: the pressure–time diagram of a supersonic object propagating shows how the transition induced by a shock wave is analogous to a dynamic phase transition.
When an object (or disturbance) moves faster than the information can propagate into the surrounding fluid, then the fluid near the disturbance cannot react or "get out of the way" before the disturbance arrives. In a shock wave the properties of the fluid (density, pressure, temperature, flow velocity, Mach number) change almost instantaneously. Measurements of the thickness of shock waves in air have resulted in values around 200 nm (about 10−5 in), which is on the same order of magnitude as the mean free path of gas molecules. In reference to the continuum, this implies the shock wave can be treated as either a line or a plane if the flow field is two-dimensional or three-dimensional, respectively.
Shock waves are formed when a pressure front moves at supersonic speeds and pushes on the surrounding air. At the region where this occurs, sound waves travelling against the flow reach a point where they cannot travel any further upstream and the pressure progressively builds in that region; a high-pressure shock wave rapidly forms.
Shock waves are not conventional sound waves; a shock wave takes the form of a very sharp change in the gas properties. Shock waves in air are heard as a loud "crack" or "snap" noise. Over longer distances, a shock wave can change from a nonlinear wave into a linear wave, degenerating into a conventional sound wave as it heats the air and loses energy. The sound wave is heard as the familiar "thud" or "thump" of a sonic boom, commonly created by the supersonic flight of aircraft.
The shock wave is one of several different ways in which a gas in a supersonic flow can be compressed. Some other methods are isentropic compressions, including Prandtl–Meyer compressions. The method of compression of a gas results in different temperatures and densities for a given pressure ratio which can be analytically calculated for a non-reacting gas. A shock wave compression results in a loss of total pressure, meaning that it is a less efficient method of compressing gases for some purposes, for instance in the intake of a scramjet. The appearance of pressure-drag on supersonic aircraft is mostly due to the effect of shock compression on the flow.
Normal shocks
In elementary fluid mechanics utilizing ideal gases, a shock wave is treated as a discontinuity where entropy increases abruptly as the shock passes. Since no fluid flow is discontinuous, a control volume is established around the shock wave, with the control surfaces that bound this volume parallel to the shock wave (with one surface on the pre-shock side of the fluid medium and one on the post-shock side). The two surfaces are separated by a very small depth such that the shock itself is entirely contained between them. At such control surfaces, momentum, mass flux and energy are constant; within combustion, detonations can be modelled as heat introduction across a shock wave. It is assumed the system is adiabatic (no heat exits or enters the system) and no work is being done. The Rankine–Hugoniot conditions arise from these considerations.
Taking into account the established assumptions, in a system where the downstream properties are becoming subsonic: the upstream and downstream flow properties of the fluid are considered isentropic. Since the total amount of energy within the system is constant, the stagnation enthalpy remains constant over both regions. However, entropy is increasing; this must be accounted for by a drop in stagnation pressure of the downstream fluid.
Other shocks
Oblique shocks
When analyzing shock waves in a flow field, which are still attached to the body, the shock wave which is deviating at some arbitrary angle from the flow direction is termed oblique shock. These shocks require a component vector analysis of the flow; doing so allows for the treatment of the flow in an orthogonal direction to the oblique shock as a normal shock.
Bow shocks
When an oblique shock is likely to form at an angle which cannot remain on the surface, a nonlinear phenomenon arises where the shock wave will form a continuous pattern around the body. These are termed bow shocks. In these cases, the 1d flow model is not valid and further analysis is needed to predict the pressure forces which are exerted on the surface.
Shock waves due to nonlinear steepening
Shock waves can form due to steepening of ordinary waves. The best-known example of this phenomenon is ocean waves that form breakers on the shore. In shallow water, the speed of surface waves is dependent on the depth of the water. An incoming ocean wave has a slightly higher wave speed near the crest of each wave than near the troughs between waves, because the wave height is not infinitesimal compared to the depth of the water. The crests overtake the troughs until the leading edge of the wave forms a vertical face and spills over to form a turbulent shock (a breaker) that dissipates the wave's energy as sound and heat.
Similar phenomena affect strong sound waves in gas or plasma, due to the dependence of the sound speed on temperature and pressure. Strong waves heat the medium near each pressure front, due to adiabatic compression of the air itself, so that high pressure fronts outrun the corresponding pressure troughs. There is a theory that the sound pressure levels in brass instruments such as the trombone become high enough for steepening to occur, forming an essential part of the bright timbre of the instruments. While shock formation by this process does not normally happen to unenclosed sound waves in Earth's atmosphere, it is thought to be one mechanism by which the solar chromosphere and corona are heated, via waves that propagate up from the solar interior.
Analogies
A shock wave may be described as the furthest point upstream of a moving object which "knows" about the approach of the object. In this description, the shock wave position is defined as the boundary between the zone having no information about the shock-driving event and the zone aware of the shock-driving event, analogous with the light cone described in the theory of special relativity.
To produce a shock wave, an object in a given medium (such as air or water) must travel faster than the local speed of sound. In the case of an aircraft travelling at high subsonic speed, regions of air around the aircraft may be travelling at exactly the speed of sound, so that the sound waves leaving the aircraft pile up on one another, similar to a traffic jam on a motorway. When a shock wave forms, the local air pressure increases and then spreads out sideways. Because of this amplification effect, a shock wave can be very intense, more like an explosion when heard at a distance (not coincidentally, since explosions create shock waves).
Analogous phenomena are known outside fluid mechanics. For example, charged particles accelerated beyond the speed of light in a refractive medium (such as water, where the speed of light is less than that in a vacuum) create visible shock effects, a phenomenon known as Cherenkov radiation.
Phenomenon types
Below are a number of examples of shock waves, broadly grouped with similar shock phenomena:
Moving shock
Usually consists of a shock wave propagating into a stationary medium
In this case, the gas ahead of the shock is stationary (in the laboratory frame) and the gas behind the shock can be supersonic in the laboratory frame. The shock propagates with a wavefront which is normal (at right angles) to the direction of flow. The speed of the shock is a function of the original pressure ratio between the two bodies of gas.
Moving shocks are usually generated by the interaction of two bodies of gas at different pressure, with a shock wave propagating into the lower pressure gas and an expansion wave propagating into the higher pressure gas.
Examples: Balloon bursting, shock tube, shock wave from explosion.
Detonation wave
A detonation wave is essentially a shock supported by a trailing exothermic reaction. It involves a wave travelling through a highly combustible or chemically unstable medium, such as an oxygen-methane mixture or a high explosive. The chemical reaction of the medium occurs following the shock wave, and the chemical energy of the reaction drives the wave forward.
A detonation wave follows slightly different rules from an ordinary shock since it is driven by the chemical reaction occurring behind the shock wavefront. In the simplest theory for detonations, an unsupported, self-propagating detonation wave proceeds at the Chapman–Jouguet flow velocity. A detonation will also cause a shock to propagate into the surrounding air due to the overpressure induced by the explosion.
When a shock wave is created by high explosives such as TNT (which has a detonation velocity of 6,900 m/s), it will always travel at high, supersonic velocity from its point of origin.
Bow shock (detached shock)
These shocks are curved and form a small distance in front of the body. Directly in front of the body, they stand at 90 degrees to the oncoming flow and then curve around the body. Detached shocks allow the same type of analytic calculations as for the attached shock, for the flow near the shock. They are a topic of continuing interest, because the rules governing the shock's distance ahead of the blunt body are complicated and are a function of the body's shape. Additionally, the shock standoff distance varies drastically with the temperature for a non-ideal gas, causing large differences in the heat transfer to the thermal protection system of the vehicle. See the extended discussion on this topic at atmospheric reentry. These follow the "strong-shock" solutions of the analytic equations, meaning that for some oblique shocks very close to the deflection angle limit, the downstream Mach number is subsonic. | Physical sciences | Waves | Physics |
146290 | https://en.wikipedia.org/wiki/Medical%20emergency | Medical emergency | A medical emergency is an acute injury or illness that poses an immediate risk to a person's life or long-term health, sometimes referred to as a situation risking "life or limb". These emergencies may require assistance from another, qualified person, as some of these emergencies, such as cardiovascular (heart), respiratory, and gastrointestinal cannot be dealt with by the victim themselves. Dependent on the severity of the emergency, and the quality of any treatment given, it may require the involvement of multiple levels of care, from first aiders through emergency medical technicians, paramedics, emergency physicians and anesthesiologists.
Any response to an emergency medical situation will depend strongly on the situation, the patient involved, and availability of resources to help them. It will also vary depending on whether the emergency occurs whilst in hospital under medical care, or outside medical care (for instance, in the street or alone at home).
Response
Summoning emergency services
For emergencies starting outside medical care, a key component of providing proper care is to summon the emergency medical services (usually an ambulance), by calling for help using the appropriate local emergency telephone number, such as 999, 911, 111, 112 or 000. After determining that the incident is a medical emergency (as opposed to, for example, a police call), the emergency dispatchers will generally run through a questioning system such as AMPDS in order to assess the priority level of the call, along with the caller's name and location.
First aid and assisting emergency services
Those who are trained to perform first aid can act within the bounds of the knowledge they have, whilst awaiting the next level of definitive care.
Those who are not able to perform first aid can also assist by remaining calm and staying with the injured or ill person. A common complaint of emergency service personnel is the propensity of people to crowd around the scene of a victim, as it is generally unhelpful, making the patient more stressed, and obstructing the smooth working of the emergency services. If possible, first responders should designate a specific person to ensure that the emergency services are called. Another bystander should be sent to wait for their arrival and direct them to the proper location. Additional bystanders can be helpful in ensuring that crowds are moved away from the ill or injured patient, allowing the responder adequate space to work.
Legal protections for responders
To prevent the delay of life-saving aid from bystanders, many states of the USA have "Good Samaritan laws" which protect civilian responders who choose to assist in an emergency. In many situations, the general public may delay giving care due to fear of liability should they accidentally cause harm. Good Samaritan laws often protect responders who act within the scope of their knowledge and training, as a "reasonable person" in the same situation would act.
The concept of implied consent can protect first responders in emergency situations. A first responder may not legally touch a patient without the patient's consent. However, consent may be either expressed or implied:
If a patient is able to make decisions, they must give expressed, informed consent before aid is given.
However, if a patient is too injured or ill to make decisions – for example, if they are unconscious, have an altered mental status, or cannot communicate - implied consent applies. Implied consent means that treatment can be given, because it is assumed that the patient would want that care.
Usually, once care has begun, a first responder or first aid provider may not leave the patient or terminate care until a responder of equal or higher training (such as an emergency medical technician) assumes care. This can constitute abandonment of the patient and may subject the responder to legal liability. Care must be continued until the patient is transferred to a higher level of care; the situation becomes too unsafe to continue; or the responder is physically unable to continue due to exhaustion or hazards.
Unless the situation is particularly hazardous and is likely to further endanger the patient, evacuating an injured victim requires special skills, and should be left to the professionals of the emergency medical and fire service.
The chain of survival
During a medical emergency in which a patient is no longer breathing and does not have a pulse, survival is predicated on adherence to the chain of survival, which has four components:
Early access to emergency care
Early cardiopulmonary resuscitation (CPR)
Early defibrillation
Early advanced life support (ALS)
Clinical response
Within hospital settings, an adequate staff is generally present to deal with the average emergency situation. Emergency medicine physicians and anaesthesiologists have training to deal with most medical emergencies, and maintain CPR and Advanced Cardiac Life Support (ACLS) certifications. In disasters or complex emergencies, most hospitals have protocols to summon on-site and off-site staff rapidly.
Both emergency department and inpatient medical emergencies follow the basic protocol of Advanced Cardiac Life Support. Irrespective of the nature of the emergency, adequate blood pressure and oxygenation are required before the cause of the emergency can be eliminated. Possible exceptions include the clamping of arteries in severe hemorrhage.
Non-trauma emergencies
While the golden hour is a trauma treatment concept, two emergency medical conditions have well-documented time-critical treatment considerations: stroke and myocardial infarction (heart attack). In the case of stroke, there is a window of three hours within which the benefit of thrombolytic drugs outweighs the risk of major bleeding. In the case of a heart attack, rapid stabilization of fatal arrhythmias can prevent sudden cardiac arrest. In addition, there is a direct relationship between time-to-treatment and the success of reperfusion (restoration of blood flow to the heart), including a time-dependent reduction in the mortality and morbidity.
| Biology and health sciences | General concepts | Health |
146311 | https://en.wikipedia.org/wiki/Shock%20%28circulatory%29 | Shock (circulatory) | Shock is the state of insufficient blood flow to the tissues of the body as a result of problems with the circulatory system. Initial symptoms of shock may include weakness, tachycardia, hyperventilation, sweating, anxiety, and increased thirst. This may be followed by confusion, unconsciousness, or cardiac arrest, as complications worsen.
Shock is divided into four main types based on the underlying cause: hypovolemic, cardiogenic, obstructive, and distributive shock. Hypovolemic shock, also known as low volume shock, may be from bleeding, diarrhea, or vomiting. Cardiogenic shock may be due to a heart attack or cardiac contusion. Obstructive shock may be due to cardiac tamponade or a tension pneumothorax. Distributive shock may be due to sepsis, anaphylaxis, injury to the upper spinal cord, or certain overdoses.
The diagnosis is generally based on a combination of symptoms, physical examination, and laboratory tests. A decreased pulse pressure (systolic blood pressure minus diastolic blood pressure) or a fast heart rate raises concerns.
Shock is a medical emergency and requires urgent medical care. If shock is suspected, emergency help should be called immediately. While waiting for medical care, the individual should be, if safe, laid down (except in cases of suspected head or back injuries). The legs should be raised if possible, and the person should be kept warm. If the person is unresponsive, breathing should be monitored and CPR may need to be performed.
Signs and symptoms
The presentation of shock is variable, with some people having only minimal symptoms such as confusion and weakness. While the general signs for all types of shock are low blood pressure, decreased urine output, and confusion, these may not always be present. While a fast heart rate is common, in those on β-blockers, those who are athletic, and in 30% of cases of those with shock due to intra abdominal bleeding, heart rate may be normal or slow. Specific subtypes of shock may have additional symptoms.
Dry mucous membrane, reduced skin turgor, prolonged capillary refill time, weak peripheral pulses, and cold extremities can be early signs of shock.
Low volume
Hypovolemic shock is the most common type of shock and is caused by insufficient circulating volume. The most common cause of hypovolemic shock is hemorrhage (internal or external); however, vomiting and diarrhea are more common causes in children. Other causes include burns, as well as excess urine loss due to diabetic ketoacidosis and diabetes insipidus.
Signs and symptoms of hypovolemic shock include:
A rapid, weak, thready pulse due to decreased blood flow combined with tachycardia
Cool skin due to vasoconstriction and stimulation of vasoconstriction
Rapid and shallow breathing due to sympathetic nervous system stimulation and acidosis
Hypothermia due to decreased perfusion and evaporation of sweat
Thirst and dry mouth, due to fluid depletion
Cold and mottled skin (livedo reticularis), especially extremities, due to insufficient perfusion of the skin
The severity of hemorrhagic shock can be graded on a 1–4 scale on the physical signs. The shock index (heart rate divided by systolic blood pressure) is a stronger predictor of the impact of blood loss than heart rate and blood pressure alone. This relationship has not been well established in pregnancy-related bleeding.
Cardiogenic
Cardiogenic shock is caused by the failure of the heart to pump effectively. This can be due to damage to the heart muscle, most often from a large myocardial infarction. Other causes of cardiogenic shock include dysrhythmias, cardiomyopathy/myocarditis, congestive heart failure (CHF), myocardial contusion, or valvular heart disease problems.
Symptoms of cardiogenic shock include:
Distended jugular veins due to increased jugular venous pressure
Weak or absent pulse
Abnormal heart rhythms, often a fast heart rate
Pulsus paradoxus in case of tamponade
Reduced blood pressure
Shortness of breath, due to pulmonary congestion
Obstructive
Obstructive shock is a form of shock associated with physical obstruction of the great vessels of the systemic or pulmonary circulation. Several conditions can result in this form of shock.
Cardiac tamponade, in which fluid in the pericardium prevents inflow of blood into the heart (venous return).
Constrictive pericarditis, in which the pericardium shrinks and hardens, is similar in presentation.
Tension pneumothorax; Through increased intrathoracic pressure, venous return is impeded.
Pulmonary embolism is thromboembolism of the lungs, hindering oxygenation and return of blood to the heart.
Aortic stenosis hinders circulation by obstructing the cardiac output.
Hypertrophic sub-aortic stenosis is overly thick ventricular muscle that dynamically occludes the ventricular outflow tract.
Abdominal compartment syndrome defined as an increase in intra-abdominal pressure to > 20 mmHg with organ dysfunction. Increased intra-abdominal pressure can result from sepsis and severe abdominal trauma. This increased pressure reduces venous return, thereby reducing lung-heart function, resulting in signs and symptoms of shock.
Many of the signs of obstructive shock are similar to cardiogenic shock, although treatments differ. Symptoms of obstructive shock include:
Abnormal heart rhythms, often a fast heart rate.
Reduced blood pressure.
Cool, clammy, mottled skin, often due to low blood pressure and vasoconstriction.
Decreased urine output.
Distributive
Distributive shock is low blood pressure due to a dilation of blood vessels within the body. This can be caused by systemic infection (septic shock), a severe allergic reaction (anaphylaxis), or spinal cord injury (neurogenic shock).
Septic shock is the most common cause of distributive shock. It is caused by an overwhelming systemic infection resulting in vasodilation leading to hypotension. Septic shock can be caused by Gram negative bacteria such as (among others) Escherichia coli, Proteus species, Klebsiella pneumoniae (which have an endotoxin on their surface that produces adverse biochemical, immunological and occasionally neurological effects which are harmful to the body), other Gram-positive cocci, such as pneumococci and streptococci, and certain fungi as well as Gram-positive bacterial toxins. Septic shock also includes some elements of cardiogenic shock. In 1992, the ACCP/SCCM Consensus Conference Committee defined septic shock: " ... sepsis-induced hypotension (systolic blood pressure < 90 mmHg or a reduction of 40 mmHg from baseline) despite adequate fluid resuscitation along with the presence of perfusion abnormalities that may include, but are not limited to: lactic acidosis, oliguria, or an acute alteration in mental status. Patients who are receiving inotropic or vasopressor agents may have a normalized blood pressure at the time that perfusion abnormalities are identified. The pathophysiology behind septic shock is as follows: 1) Systemic leukocyte adhesion to endothelial cells 2) Reduced contractility of the heart 3) Activation of the coagulation pathways, resulting in disseminated intravascular coagulation 4). Increased levels of neutrophils
The main manifestations of septic shock are due to the massive release of histamine which causes intense dilation of the blood vessels. People with septic shock will also likely be positive for SIRS criteria. The most generally accepted treatment for these patients is early recognition of symptoms, and early administration of broad spectrum and organism specific antibiotics.
Signs of septic shock include:
Abnormal heart rhythms, often a fast heart rate
Reduced blood pressure
Decreased urine output
Altered mental status
Anaphylactic shock is caused by a severe anaphylactic reaction to an allergen, antigen, drug, or foreign protein causing the release of histamine which causes widespread vasodilation, leading to hypotension and increased capillary permeability. Signs typically occur after exposure to an allergen and may include:
Skin changes, such as hives, itching, flushing, and swelling.
Wheezing and shortness of breath.
Abdominal pain, diarrhea, and vomiting.
Lightheadedness, confusion, headaches, loss of consciousness.
High spinal injuries may cause neurogenic shock, which is commonly classified as a subset of distributive shock. The classic symptoms include a slow heart rate due to loss of cardiac sympathetic tone and warm skin due to dilation of the peripheral blood vessels. (This term can be confused with spinal shock which is a recoverable loss of function of the spinal cord after injury and does not refer to the hemodynamic instability.)
Endocrine
Although not officially classified as a subcategory of shock, many endocrinological disturbances in their severe form can result in shock.
Hypothyroidism (can be considered a form of cardiogenic shock) in people who are critically ill patients reduces cardiac output and can lead to hypotension and respiratory insufficiency.
Thyrotoxicosis (cardiogenic shock) may induce a reversible cardiomyopathy.
Acute adrenal insufficiency (distributive shock) is frequently the result of discontinuing corticosteroid treatment without tapering the dosage. However, surgery and intercurrent disease in patients on corticosteroid therapy without adjusting the dosage to accommodate for increased requirements may also result in this condition.
Relative adrenal insufficiency (distributive shock) in critically ill patients where present hormone levels are insufficient to meet the higher demands.
Cause
Shock is a common end point of many medical conditions. Shock triggered by a serious allergic reaction is known as anaphylactic shock, shock triggered by severe dehydration or blood loss is known as hypovolemic shock, shock caused by sepsis is known as septic shock, etc. Shock itself is a life-threatening condition as a result of compromised body circulation. It can be divided into four main types based on the underlying cause: hypovolemic, distributive, cardiogenic, and obstructive. A few additional classifications are occasionally used, such as endocrinologic shock.
Pathophysiology
Shock is a complex and continuous condition, and there is no sudden transition from one stage to the next. At a cellular level, shock is the process of oxygen demand becoming greater than oxygen supply.
One of the key dangers of shock is that it progresses by a positive feedback loop. Poor blood supply leads to cellular damage, which results in an inflammatory response to increase blood flow to the affected area. Normally, this causes the blood supply level to match with tissue demand for nutrients. However, if there is enough increased demand in some areas, it can deprive other areas of sufficient supply, which then start demanding more. This then leads to an ever escalating cascade.
As such, shock is a runaway condition of homeostatic failure, where the usual corrective mechanisms relating to oxygenation of the body no longer function in a stable way. When it occurs, immediate treatment is critical in order to return an individual's metabolism into a stable, self-correcting trajectory. Otherwise the condition can become increasingly difficult to correct, surprisingly quickly, and then progress to a fatal outcome. In the particular case of anaphylactic shock, progression to death might take just a few minutes.
Initial
During the Initial stage (Stage 1), the state of hypoperfusion causes hypoxia. Due to the lack of oxygen, the cells perform lactic acid fermentation. Since oxygen, the terminal electron acceptor in the electron transport chain, is not abundant, this slows down entry of pyruvate into the Krebs cycle, resulting in its accumulation. The accumulating pyruvate is converted to lactate (lactic acid) by lactate dehydrogenase. The accumulating lactate causes lactic acidosis.
Compensatory
The Compensatory stage (Stage 2) is characterised by the body employing physiological mechanisms, including neural, hormonal, and bio-chemical mechanisms, in an attempt to reverse the condition. As a result of the acidosis, the person will begin to hyperventilate in order to rid the body of carbon dioxide (CO2) since it indirectly acts to acidify the blood; the body attempts to return to acid–base homeostasis by removing that acidifying agent. The baroreceptors in the arteries detect the hypotension resulting from large amounts of blood being redirected to distant tissues, and cause the release of epinephrine and norepinephrine. Norepinephrine causes predominately vasoconstriction with a mild increase in heart rate, whereas epinephrine predominately causes an increase in heart rate with a small effect on the vascular tone; the combined effect results in an increase in blood pressure. The renin–angiotensin axis is activated, and arginine vasopressin (anti-diuretic hormone) is released to conserve fluid by reducing its excretion via the renal system. These hormones cause the vasoconstriction of the kidneys, gastrointestinal tract, and other organs to divert blood to the heart, lungs and brain. The lack of blood to the renal system causes the characteristic low urine production. However, the effects of the renin–angiotensin axis take time and are of little importance to the immediate homeostatic mediation of shock.
Progressive/decompensated
The Progressive stage (stage 3) results if the underlying cause of the shock is not successfully treated. During this stage, compensatory mechanisms begin to fail. Due to the decreased perfusion of the cells in the body, sodium ions build up within the intracellular space while potassium ions leak out. Due to lack of oxygen, cellular respiration diminishes and anaerobic metabolism predominates. As anaerobic metabolism continues, the arteriolar smooth muscle and precapillary sphincters relax such that blood remains in the capillaries. Due to this, the hydrostatic pressure will increase and, combined with histamine release, will lead to leakage of fluid and protein into the surrounding tissues. As this fluid is lost, the blood concentration and viscosity increase, causing sludging of the micro-circulation. The prolonged vasoconstriction will also cause the vital organs to be compromised due to reduced perfusion. If the bowel becomes sufficiently ischemic, bacteria may enter the blood stream, resulting in the increased complication of endotoxic shock.
Refractory
At Refractory stage (stage 4), the vital organs have failed and the shock can no longer be reversed. Brain damage and cell death are occurring, and death will occur imminently. One of the primary reasons that shock is irreversible at this point is that much of the cellular ATP (the basic energy source for cells) has been degraded into adenosine in the absence of oxygen as an electron receptor in the mitochondrial matrix. Adenosine easily perfuses out of cellular membranes into extracellular fluid, furthering capillary vasodilation, and then is transformed into uric acid. Because cells can only produce adenosine at a rate of about 2% of the cell's total need per hour, even restoring oxygen is futile at this point because there is no adenosine to phosphorylate into ATP.
Diagnosis
The diagnosis of shock is commonly based on a combination of symptoms, physical examination, and laboratory tests. Many signs and symptoms are not sensitive or specific for shock, thus many clinical decision-making tools have been developed to identify shock at an early stage. A high degree of suspicion is necessary for the proper diagnosis of shock.
Shock is, hemodynamically speaking, inadequate blood flow or cardiac output, Unfortunately, the measurement of cardiac output requires an invasive catheter, such as a pulmonary artery catheter. Mixed venous oxygen saturation (SmvO2) is one of the methods of calculating cardiac output with a pulmonary artery catheter. Central venous oxygen saturation (ScvO2) as measured via a central line correlates well with SmvO2 and is easier to acquire.
Tissue oxygenation
is critically dependent on blood flow. When the oxygenation of tissues is compromised anaerobic metabolism will begin and lactic acid will be produced.
Management
Treatment of shock is based on the likely underlying cause. An open airway and sufficient breathing should be established. Any ongoing bleeding should be stopped, which may require surgery or embolization. Intravenous fluid, such as Ringer's lactate or packed red blood cells, is often given. Efforts to maintain a normal body temperature are also important. Vasopressors may be useful in certain cases. Shock is both common and has a high risk of death. In the United States about 1.2 million people present to the emergency room each year with shock and their risk of death is between 20 and 50%.
The best evidence exists for the treatment of septic shock in adults. However, the pathophysiology of shock in children appears to be similar so treatment methodologies have been extrapolated to children. Management may include securing the airway via intubation if necessary to decrease the work of breathing and for guarding against respiratory arrest. Oxygen supplementation, intravenous fluids, passive leg raising (not Trendelenburg position) should be started and blood transfusions added if blood loss is severe. In select cases, compression devices like non-pneumatic anti-shock garments (or the deprecated military anti-shock trousers) can be used to prevent further blood loss and concentrate fluid in the body's head and core. It is important to keep the person warm to avoid hypothermia as well as adequately manage pain and anxiety as these can increase oxygen consumption. Negative impact by shock is reversible if it's recognized and treated early in time.
Fluids
Aggressive intravenous fluids are recommended in most types of shock (e.g. 1–2 liter normal saline bolus over 10 minutes or 20 mL/kg in a child) which is usually instituted as the person is being further evaluated. Colloids and crystalloids appear to be equally effective with respect to outcomes., Balanced crystalloids and normal saline also appear to be equally effective in critically ill patients. If the person remains in shock after initial resuscitation, packed red blood cells should be administered to keep the hemoglobin greater than 100 g/L.
For those with hemorrhagic shock, the current evidence supports limiting the use of fluids for penetrating thorax and abdominal injuries allowing mild hypotension to persist (known as permissive hypotension). Targets include a mean arterial pressure of 60 mmHg, a systolic blood pressure of 70–90 mmHg, or until the patient has adequate mentation and peripheral pulses. Hypertonic fluid may also be an option in this group.
Medications
Vasopressors may be used if blood pressure does not improve with fluids. Common vasopressors used in shock include: norepinephrine, phenylephrine, dopamine, and dobutamine.
There is no evidence of substantial benefit of one vasopressor over another; however, using dopamine leads to an increased risk of arrhythmia when compared with norepinephrine. Vasopressors have not been found to improve outcomes when used for hemorrhagic shock from trauma but may be of use in neurogenic shock. Activated protein C (Xigris), while once aggressively promoted for the management of septic shock, has been found not to improve survival and is associated with a number of complications. Activated protein C was withdrawn from the market in 2011, and clinical trials were discontinued. The use of sodium bicarbonate is controversial as it has not been shown to improve outcomes. If used at all it should only be considered if the blood pH is less than 7.0.
People with anaphylactic shock are commonly treated with epinephrine. Antihistamines, such as Benadryl (diphenhydramine) or ranitidine are also commonly administered. Albuterol, normal saline, and steroids are also commonly given.
Mechanical support
Intra-aortic balloon pump (IABP) – a device inserted into the aorta that mechanically raises the blood pressure. Use of Intra-aortic balloon pumps is not recommended in cardiogenic shock.
Ventricular assist device (VAD) – A mechanical pump that helps pump blood throughout the body. Commonly used in short term cases of refractory primary cardiogenic shock.
Artificial heart (TAH)
Extracorporeal membrane oxygenation (ECMO) – an external device that completely replaces the work of the heart.
Treatment goals
The goal of treatment is to achieve a urine output of greater than 0.5 mL/kg/h, a central venous pressure of 8–12 mmHg and a mean arterial pressure of 65–95 mmHg. In trauma the goal is to stop the bleeding which in many cases requires surgical interventions. A good urine output indicates that the kidneys are getting enough blood flow.
Epidemiology
Septic shock (a form of distributive shock) is the most common form of shock. Shock from blood loss occurs in about 1–2% of trauma cases. Overall, up to one-third of people admitted to the intensive care unit (ICU) are in circulatory shock. Of these, cardiogenic shock accounts for approximately 20%, hypovolemic about 20%, and septic shock about 60% of cases.
Prognosis
The prognosis of shock depends on the underlying cause and the nature and extent of concurrent problems. Low volume, anaphylactic, and neurogenic shock are readily treatable and respond well to medical therapy. Septic shock, especially septic shock where treatment is delayed or the antimicrobial drugs are ineffective, however has a mortality rate between 30% and 80%; cardiogenic shock has a mortality rate of up to 70% to 90%, though quick treatment with vasopressors and inotropic drugs, cardiac surgery, and the use of assistive devices can lower the mortality.
History
There is no evidence of the word shock being used in its modern-day form prior to 1743. However, there is evidence that Hippocrates used the word exemia to signify a state of being "drained of blood". Shock or "choc" was first described in a trauma victim in the English translation of Henri-François LeDran's 1740 text, Traité ou Reflexions Tire'es de la Pratique sur les Playes d'armes à feu (A treatise, or reflections, drawn from practice on gun-shot wounds.) In this text he describes "choc" as a reaction to the sudden impact of a missile. However, the first English writer to use the word shock in its modern-day connotation was James Latta, in 1795.
Prior to World War I, there were several competing hypotheses behind the pathophysiology of shock. Of the various theories, the most well regarded was a theory penned by George W. Crile who suggested in his 1899 monograph, "An Experimental Research into Surgical Shock", that shock was quintessentially defined as a state of circulatory collapse (vasodilation) due to excessive nervous stimulation. Other competing theories around the turn of the century included one penned by Malcom in 1907, in which the assertion was that prolonged vasoconstriction led to the pathophysiological signs and symptoms of shock. In the following World War I, research concerning shock resulted in experiments by Walter B. Cannon of Harvard and William M. Bayliss of London in 1919 that showed that an increase in permeability of the capillaries in response to trauma or toxins was responsible for many clinical manifestations of shock. In 1972 Hinshaw and Cox suggested the classification system for shock which is still used today.
| Biology and health sciences | Injury | null |
146315 | https://en.wikipedia.org/wiki/Bleeding | Bleeding | Bleeding, hemorrhage, haemorrhage or blood loss is blood escaping from the circulatory system from damaged blood vessels. Bleeding can occur internally, or externally either through a natural opening such as the mouth, nose, ear, urethra, vagina or anus, or through a puncture in the skin.
Hypovolemia is a massive decrease in blood volume, and death by excessive loss of blood is referred to as exsanguination. Typically, a healthy person can endure a loss of 10–15% of the total blood volume without serious medical difficulties (by comparison, blood donation typically takes 8–10% of the donor's blood volume). The stopping or controlling of bleeding is called hemostasis and is an important part of both first aid and surgery.
Types
Upper head
Intracranial hemorrhage — bleeding in the skull.
Cerebral hemorrhage — a type of intracranial hemorrhage, bleeding within the brain tissue itself.
Intracerebral hemorrhage — bleeding in the brain caused by the rupture of a blood vessel within the head. | Biology and health sciences | Injury | null |
146317 | https://en.wikipedia.org/wiki/Ambulance | Ambulance | An ambulance is a medically-equipped vehicle used to transport patients to treatment facilities, such as hospitals. Typically, out-of-hospital medical care is provided to the patient during the transport. Ambulances are used to respond to medical emergencies by emergency medical services (EMS), and can rapidly transport paramedics and other first responders, carry equipment for administering emergency care, and transport patients to hospital or other definitive care. Most ambulances use a design based on vans or pickup trucks, though others take the form of motorcycles, buses, hearses, aircraft and boats.
Ambulances are generally considered emergency vehicles authorized to be equipped with emergency lights and sirens. Generally, vehicles count as an ambulance if they can transport patients. However, it varies by jurisdiction as to whether a non-emergency patient transport vehicle (also called an ambulette) is counted as an ambulance. These vehicles are not usually (although there are exceptions) equipped with life-support equipment, and are usually crewed by staff with fewer qualifications than the crew of emergency ambulances. Conversely, EMS agencies may also have nontransporting EMS vehicles that cannot transport patients.
The term ambulance comes from the Latin word as meaning 'to walk or move about' which is a reference to early medical care where patients were moved by lifting or wheeling. The word originally meant a moving hospital, which follows an army in its movements. Ambulances ( in Spanish) were first used for emergency transport in 1487 by the Spanish forces during the siege of Málaga by the Catholic Monarchs against the Emirate of Granada. During the American Civil War vehicles for conveying the wounded off the field of battle were called ambulance wagons. Field hospitals were still called ambulances during the Franco-Prussian War of 1870 and in the Serbo-Turkish war of 1876 even though the wagons were first referred to as ambulances about 1854 during the Crimean War.
History
The history of the ambulance begins in ancient times, with the use of carts to transport incurable patients by force. Ambulances were first used for emergency transport in 1487 by the Spanish, and civilian variants were put into operation during the 1830s. Advances in technology throughout the 19th and 20th centuries led to modern self-powered ambulances.
Functional types
Ambulances can be grouped into types depending on whether or not they transport patients, and under what conditions. In some cases, ambulances may fulfill more than one function (such as combining emergency ambulance care with patient transport:
Emergency ambulance – The most common type of ambulance, which provides care to patients with an acute illness or injury. These can be road-going vans, boats, helicopters, fixed-wing aircraft (known as air ambulances), or even converted vehicles such as golf carts.
Patient transport ambulance – A vehicle, which has the job of transporting patients to, from or between places of medical treatment, such as hospital or dialysis center, for non-urgent care. These can be vans, buses, or other vehicles.
Ambulance bus – A large ambulance, usually based on a bus chassis, that can evacuate and transport a large number of patients. They are usually used in emergencies such as mass casualty incidents.
Charity ambulance – A special type of patient transport ambulance is provided by a charity for the purpose of taking sick children or adults on trips or vacations away from hospitals, hospices, or care homes where they are in long-term care. Examples include the United Kingdom's "Jumbulance" project. These are usually based on a bus.
Bariatric ambulance – A special type of patient transport ambulance designed for extremely obese patients equipped with the appropriate tools to move and manage these patients.
Rapid organ recovery ambulance – A special ambulance used to collect the bodies of people who have died to preserve their organs. In 2008, New York City launched a pilot program deploying one such ambulance with a $1.5 million three-year grant.
Psychiatric ambulance – A special ambulance dedicated to treat psychiatric emergencies. The idea was first tested as "Psykebilen" ("The Psych ambo") in Bergen, Norway in 2005, and was soon adopted by other cities in Norway and Sweden. Tests from the time showed that an ambulance service with personnel specially trained in psychiatric treatment was highly effective, and reduced the use of force when treating patients in psychiatric crises.
Vehicle types
Ambulances can be based on many types of vehicle although emergency and disaster conditions may lead to other vehicles serving as makeshift ambulances:
Van or pickup truck – A typical general-purpose ambulance is based on either the chassis of a van ("vanbulance") or a light-duty truck. This chassis is then modified to the designs and specifications of the purchaser. Vans may either retain their original body and be upfitted inside, or may be based on a chassis without the original body with a modular box body fitted instead. Those based on pickup trucks almost always have modular bodies. Those vehicles intended for especially intensive care or require a large amount of equipment to be carried may be based on medium-duty trucks.
Car – Used either as a fly-car for rapid response or to transport patients who can sit, these are standard car models adapted to the requirements of the service using them. Some cars are capable of taking a stretcher with a recumbent patient, but this often requires the removal of the front passenger seat, or the use of a particularly long car. This was often the case with early ambulances, which were converted (or even serving) hearses, as these were some of the few vehicles able to accept a human body in a supine position. Some operators use modular-body transport ambulances based on the chassis of a minivan or station wagon.
Motorcycle and motor scooter – In urban areas, these may be used for rapid response in an emergency as they can travel through heavy traffic much faster than a car or van. Trailers or sidecars can make these patient transporting units.
Bicycle – Used for response, but usually in pedestrian-only areas where large vehicles find access difficult. Like the motorcycle ambulance, a bicycle may be connected to a trailer for patient transport, most often in the developing world.
All-terrain vehicle – Used for response off-road, especially at events or in remote areas. ATVs can be modified to carry a stretcher, and are used for tasks such as mountain rescue in inaccessible areas.
Golf cart or Neighborhood Electric Vehicle – Used for rapid response at events or on campuses. These function similarly to ATVs, with less rough terrain capability and less noise.
Bus – In some cases, buses can be used to transport multiple casualties, either for the purposes of taking patients on journeys, in the context of major incidents, or to deal with specific problems such as drunken patients in town centers. They are sometimes referred to as ambulance buses.
Helicopter – Usually used for emergency care, either in places inaccessible by road, or in areas where speed is of the essence, as they are able to travel significantly faster than a road ambulance. Helicopter and fixed-wing ambulances are discussed in greater detail at air ambulance.
Fixed-wing aircraft – These can be used for either acute emergency care in remote areas (such as in Australia, with the 'Flying Doctors'), for patient transport over long distances (e.g. a re-patriation following an illness or injury in a foreign country), or transportation between distant hospitals. Helicopter and fixed-wing ambulances are discussed in greater detail at air ambulance.
Boat – Boats can be used to serve as water ambulances, especially in island areas or in areas with a large number of canals. Some lifeboats or lifeguard vessels may fit the description of an ambulance as they are used to transport a casualty.
Train – In remote or hard-to-reach areas that are accessed primarily by railway connections, trains may be used to provide medical care as a mobile treatment facility or transport patients to better care in more accessible areas. These are generally called hospital trains.
Trailer – In some instances, a trailer, which can be towed behind a self-propelled vehicle, can be used as an ambulance. This permits flexibility in areas with minimal access to vehicles, such as on small islands.
Horse and cart – More traditional form of transport, mostly seen in developing economies.
Fire engine – Fire services (especially in North America) often train firefighters to respond to medical emergencies and most apparatuses carry at least basic medical supplies. By design, most apparatuses cannot transport patients unless they can sit in the cab. However, some fire trucks may be designed to have a large ambulance compartment behind the front of the cab, where the driver and officer's seats are located.
Vehicle type gallery
Design and construction
Ambulance design must take into account local conditions and infrastructure. Maintained roads are necessary for road-going ambulances to arrive on scene and then transport the patient to a hospital, though in rugged areas four-wheel drive or all-terrain vehicles can be used. Fuel must be available and service facilities are necessary to maintain the vehicle.
Methods of summoning (e.g. telephone) and dispatching ambulances usually rely on electronic equipment, which itself often relies on an intact power grid. Similarly, modern ambulances are equipped with two-way radios or cellular telephones to enable them to contact hospitals, either to notify the appropriate hospital of the ambulance's pending arrival, or, in cases where physicians do not form part of the ambulance's crew, to confer with a physician for medical oversight.
Ambulances often have two stages of manufacturing. The first is frequently the manufacture of light or medium truck chassis-cabs or full-size vans (or in some places, cars) such as Mercedes-Benz, Nissan, Toyota, or Ford. The second manufacturer (known as second stage manufacturer) modifies the vehicle (which is sometimes purchased incomplete, having no body or interior behind the driver's seat) and turns it into an ambulance by adding bodywork, emergency vehicle equipment, and interior fittings. This is done by one of two methods – either coachbuilding, where the modifications are started from scratch and built on to the vehicle, or using a modular system, where a pre-built 'box' is put on to the empty chassis of the ambulance, and then finished off.
Modern ambulances are typically powered by internal combustion engines, which can be powered by any conventional fuel, including diesel, gasoline or liquefied petroleum gas, depending on the preference of the operator and the availability of different options. Colder regions often use gasoline-powered engines, as diesels can be difficult to start when they are cold. Warmer regions may favor diesel engines, as they are more efficient and more durable. Diesel power is sometimes chosen due to safety concerns, after a series of fires involving gasoline-powered ambulances during the 1980s. These fires were ultimately attributed in part to gasoline's higher volatility in comparison to diesel fuel. The type of engine may be determined by the manufacturer: in the past two decades, Ford would only sell vehicles for ambulance conversion if they are diesel-powered. Beginning in 2010, Ford will sell its ambulance chassis with a gasoline engine in order to meet emissions requirements.
In the United Kingdom, the National Health Service has set a target for all ambulances to be fully electric as part of the Net Zero campaign by 2045.
Standards
Many regions have prescribed standards which ambulances should, or must, meet in order to be used for their role. These standards may have different levels which reflect the type of patient which the ambulance is expected to transport (for instance specifying a different standard for routine patient transport than high dependency), or may base standards on the size of vehicle.
For instance, in Europe, the European Committee for Standardization publishes the standard CEN 1789, which specifies minimum compliance levels across the build of ambulance, including crash resistance, equipment levels, and exterior marking. In the United States, standards for ambulance design have existed since 1976, where the standard is published by the General Services Administration and known as KKK-A-1822. This standard has been revised several times, and is currently in version 'F' change notice #13, known as KKK-A-1822F. The National Fire Protection Association has also published a design standard, NFPA 1917, which offers an alternative to KKK-A-1822F. The Commission on Accreditation of Ambulance Services (CAAS) has published its Ground Vehicle Standard for Ambulances v2.0, effective July 2019. This standard is similar to the KKK-A-1822F and NFPA 1917–2019 specifications. The decision on which of the current (3) standards to require is left up to each individual state legislature or EMS director. Some states have no specific requirement, while others specify which standard is acceptable. Others, yet, allow the end user to decide which standard to comply to.
In the United States and Canada, there are four types of ambulances: Type I, Type II, Type III, and Type IV. Type I is based on a heavy truck chassis-cab with a custom rear compartment that is often referred to as a "box" or "module", primarily used for Advanced Life Support (ALS) or Mobile Intensive Care Unit (MICU), as well as rescue work. Type II is based on a commercial heavy-duty van with few modifications except for a raised roof and a secondary air conditioning unit for the rear of the vehicle, primarily used for Basic Life Support (BLS) and transfer of patients, though they are occasionally also used for ALS and rescue. Type III is a van chassis-cab with a custom-made rear compartment, used for ALS and rescue. Type IV is for ad hoc patient transfer using smaller utility vehicles selected for maneuverability in special environments such as dense crowds at events; these are uncommon and are not subject to federal regulations in the United States.
The move towards standardisation is now reaching countries without a history of prescriptive codes, such as India, which approved its first national standard for ambulance construction in 2013.
Safety
Ambulances, like other emergency vehicles, are required to operate in most weather conditions, including those during which civilian drivers often elect to stay off the road. Also, the ambulance crew's responsibilities to their patient often preclude their use of safety devices such as seat belts. Research has shown that ambulances are more likely to be involved in motor vehicle collisions resulting in injury or death than either fire trucks or police cars. Unrestrained occupants, particularly those riding in the patient-care compartment, are particularly vulnerable. When compared to civilian vehicles of similar size, one study found that on a per-accident basis, ambulance collisions tend to involve more people, and result in more injuries. An 11-year retrospective study concluded in 2001 found that although most fatal ambulance crashes in the United States occurred during emergency runs, they typically occurred on improved, straight, dry roads, during clear weather. Furthermore, paramedics are also at risk in ambulances while helping patients, as 27 paramedics died during ambulance trips in the US between 1991 and 2006.
Equipment
In addition to the equipment directly used for the treatment of patients, ambulances may be fitted with a range of additional equipment which is used in order to facilitate patient care. This could include:
Two-way radio – One of the most important pieces of equipment in modern emergency medical services as it allows for the issuing of jobs to the ambulance, and can allow the crew to pass information back to control or to the hospital (for example a priority ASHICE message to alert the hospital of the impending arrival of a critical patient.) More recently many services worldwide have moved from traditional analog UHF/VHF sets, which can be monitored externally, to more secure digital systems, such as those working on a GSM system, such as TETRA.
Mobile data terminal – Some ambulances are fitted with mobile data terminals (or MDTs), which are connected wirelessly to a central computer, usually at the control center. These terminals can function instead of or alongside the two-way radio and can be used to pass details of jobs to the crew, and can log the time the crew was mobile to a patient, arrived, and left the scene, or fulfill any other computer-based function.
Evidence gathering CCTV – Some ambulances are now being fitted with video cameras used to record activity either inside or outside the vehicle. They may also be fitted with sound recording facilities. This can be used as a form of protection from violence against ambulance crews, or in some cases (dependent on local laws) to prove or disprove cases where a member of the crew stands accused of malpractice.
Tail lift or ramp – Ambulances can be fitted with a tail lift or ramp in order to facilitate loading a patient without having to undertake any lifting. This is especially important where the patient is obese or specialty care transports that require large, bulky equipment such as a neonatal incubator or hospital beds. There may also be equipment linked to this such as winches which are designed to pull heavy patients into the vehicle.
Trauma lighting – In addition to normal working lighting, ambulances can be fitted with special lighting (often blue or red) which is used when the patient becomes photosensitive.
Air conditioning – Ambulances are often fitted with a separate air conditioning system to serve the working area from that which serves the cab. This helps to maintain an appropriate temperature for any patients being treated but may also feature additional features such as filtering against airborne pathogens.
Data recorders – These are often placed in ambulances to record such information as speed, braking power and time, activation of active emergency warnings such as lights and sirens, as well as seat belt usage. These are often used in coordination with GPS units.
Intermediate technology
In parts of the world that lack a high level of infrastructure, ambulances are designed to meet local conditions, being built using intermediate technology. Ambulances can also be trailers, which are pulled by bicycles, motorcycles, tractors, or animals. Animal-powered ambulances can be particularly useful in regions that are subject to flooding. Motorcycles fitted with sidecars (or motorcycle ambulances) are also used, though they are subject to some of the same limitations as more traditional over-the-road ambulances. The level of care provided by these ambulances varies between merely providing transport to a medical clinic to providing on-scene and continuing care during transport.
The design of intermediate technology ambulances must take into account not only the operation and maintenance of the ambulance, but its construction as well. The robustness of the design becomes more important, as does the nature of the skills required to properly operate the vehicle. Cost-effectiveness can be a high priority.
Appearance and markings
Emergency ambulances are highly likely to be involved in hazardous situations, including incidents such as a road traffic collision, as these emergencies create people who are likely to be in need of treatment. They are required to gain access to patients as quickly as possible, and in many countries, are given dispensation from obeying certain traffic laws. For instance, they may be able to treat a red traffic light or stop sign as a yield sign ('give way'), or be permitted to break the speed limit. Generally, the priority of the response to the call will be assigned by the dispatcher, but the priority of the return will be decided by the ambulance crew based on the severity of the patient's illness or injury. Patients in significant danger to life and limb (as determined by triage) require urgent treatment by advanced medical personnel, and because of this need, emergency ambulances are often fitted with passive and active visual and/or audible warnings to alert road users.
Passive visual warnings
Passive visual warnings are usually part of the design of the vehicle, and involve the use of high contrast patterns. Older ambulances (and those in developing countries) are more likely to have their pattern painted on, whereas modern ambulances generally carry retro-reflective designs, which reflects light from car headlights or torches. Popular patterns include 'checker board' (alternate coloured squares, sometimes called 'Battenburg', named after a type of cake), chevrons (arrowheads – often pointed towards the front of the vehicle if on the side, or pointing vertically upwards on the rear) or stripes along the side (these were the first type of retro-reflective device introduced, as the original reflective material, invented by 3M, only came in tape form). In addition to retro-reflective markings, some services now have the vehicles painted in a bright (sometimes fluorescent) yellow or orange for maximum visual impact, though classic white or red are also common. Fire department-operated ambulances are often painted red to match the fire apparatuses.
Another passive marking form is the word ambulance (or local language variant) spelled out in reverse on the front of the vehicle. This enables drivers of other vehicles to more easily identify an approaching ambulance in their rear view mirrors. Ambulances may display the name of their owner or operator, and an emergency telephone number for the ambulance service.
Ambulances may also carry an emblem (either as part of the passive warning markings or not), such as a Red Cross, Red Crescent or Red Crystal (collective known as the Protective Symbols). These are symbols laid down by the Geneva Convention, and all countries signatory to it agree to restrict their use to either (1) Military Ambulances or (2) the national Red Cross or Red Crescent society. Use by any other person, organization or agency is in breach of international law. The protective symbols are designed to indicate to all people (especially combatants in the case of war) that the vehicle is neutral and is not to be fired upon, hence giving protection to the medics and their casualties, although this has not always been adhered to. In Israel, Magen David Adom, the Red Cross member organization use a red Star of David, but this does not have recognition beyond Israeli borders, where they must use the Red Crystal.
The Star of Life is widely used, and was originally designed and governed by the U.S. National Highway Traffic Safety Administration, because the Red Cross symbol is legally protected by both National and international law.
Ambulance services with historical origins such as the Order of St John, the Order of Malta Ambulance Corps and Malteser International often use the Maltese cross to identify their ambulances. This is especially important in countries such as Australia, where St. John Ambulance operate one state and one territory ambulance service, and all of Australia's other ambulance services use variations on a red Maltese cross.
Fire service operated ambulances may display the Cross of St. Florian (often incorrectly called a Maltese cross) as this cross is frequently used as a fire department logo (St. Florian being the patron saint of firefighters).
Active visual warnings
The active visual warnings are usually in the form of flashing lights. These flash in order to attract the attention of other road users as the ambulance approaches, or to provide warning to motorists approaching a stopped ambulance in a dangerous position on the road. Common colours for ambulance warning beacons are blue, red, amber, and white (clear). However the colours may vary by country and sometimes by operator.
There are several technologies in use to achieve the flashing effect. These include flashing a light bulb or LED, flashing or rotating halogen, and strobe lights, which are usually brighter than incandescent lights. Each of these can be programmed to flash singly or in groups, and can be programmed to flash in patterns (such as a left -> right pattern for use when the ambulance is parked on the left hand side of the road, indicating to other road users that they should move to the right (away from the ambulance)). Incandescent and LED lights may also be programmed to burn steadily, without flashing, which is required in some provinces.
Emergency lights may simply be mounted directly on the body, or may be housed in special fittings, such as in a lightbar or in special flush-mount designs (as seen on the Danish ambulance to the right), or may be hidden in a host light (such as a headlamp) by drilling a hole in the host light's reflector and inserting the emergency light. These hidden lights may not be apparent until they are activated. Additionally, some of the standard lights fitted to an ambulance (e.g. headlamps, tail lamps) may be programmed to flash. Flashing headlights (typically the high beams, flashed alternately) are known as a wig-wag. Additional white lights may be placed strategically around the vehicle to illuminate the area around it when it is dark, almost always at the rear for loading and unloading stretchers and often at the sides as well. In areas very far North or South where there are times of year with long periods of darkness, additional driving lights at the front are often fitted as well to increase visibility for the driver.
In order to increase safety, it is best practice to have 360° coverage with the active warnings, improving the chance of the vehicle being seen from all sides. In some countries, such as the United States, this may be mandatory. The roof, front grille, sides and rear of the body, and front fenders are common places to mount emergency lights. A certain balance must be made when deciding on the number and location of lights: too few and the ambulance may not be noticed easily, too many and it becomes a massive distraction for other road users more than it is already, increasing the risk of local accidents.
Audible warnings
In addition to visual warnings, ambulances can be fitted with audible warnings, sometimes known as sirens, which can alert people and vehicles to the presence of an ambulance before they can be seen. The first audible warnings were mechanical bells, mounted to either the front or roof of the ambulance. Most modern ambulances are now fitted with electronic sirens, producing a range of different noises which ambulance operators can use to attract more attention to themselves, particularly when proceeding through an intersection or in heavy traffic.
The speakers for modern sirens can be integral to the lightbar, or they may be hidden in or flush to the grill to reduce noise inside the ambulance that may interfere with patient care and radio communications. Ambulances can additionally be fitted with airhorn audible warnings to augment the effectiveness of the siren system, or may be fitted with extremely loud two-tone air horns as their primary siren.
A recent development is the use of the RDS system of car radios. The ambulance is fitted with a short range FM transmitter, set to RDS code 31, which interrupts the radio of all cars within range, in the manner of a traffic broadcast, but in such a way that the user of the receiving radio is unable to opt-out of the message (as with traffic broadcasts). This feature is built into every RDS radio for use in national emergency broadcast systems, but short-range units on emergency vehicles can prove an effective means of alerting traffic to their presence. It is, however, unlikely that this system could replace audible warnings, as it is unable to alert pedestrians, those not using a compatible radio or even have it turned off.
Costs
In the United States, the cost of an ambulance ride may be paid for from several sources, and this will depend on the local situation type of service being provided, by whom, and to whom.
Government-funded service – The full or the majority of the cost of transport by ambulance is borne by the local, regional, or national government (through their normal taxation).
Privately funded service – Transport by ambulance is paid for by the patient themselves, or through their insurance company. This may be at the point of care (i.e. payment or guarantee must be made before treatment or transport), although this may be an issue with critically injured patients, unable to provide such details, or via a system of billing later on.
Charity-funded service – Transport by ambulance may be provided free of charge to patients by a charity, although donations may be sought for services received.
Hospital-funded service – Hospitals may provide the ambulance transport free of charge, on the condition that patients use the hospital's services (which they may have to pay for).
Crewing
There are differing levels of qualification that the ambulance crew may hold, from holding no formal qualification to having a fully qualified doctor on board. Most ambulance services require at least two crew members to be on every ambulance (one to drive, and one to attend the patient). It may be the case that only the attendant need be qualified, and the driver might have no medical training. In some locations, an advanced life support ambulance may be crewed by one paramedic and one technician, or in countries like Australia advanced life support registered paramedics.
Common ambulance crew qualifications are:
First responder – A person who arrives first at the scene of an incident, and whose job is to provide early critical care such as cardiopulmonary resuscitation (CPR) or using an automated external defibrillator (AED). First responders may be dispatched by the ambulance service, may be passers-by, or may be dispatched to the scene from other agencies, such as the police or fire departments. They may be on duty for another agency, or volunteers who are on-call during their free time.
Ambulance driver – Some services employ staff with no medical qualification (or just basic first aid training) whose job is to simply drive the vehicle. In some emergency ambulance contexts this term is a pejorative towards personnel with higher medical training, as it implies they perform no function other than driving, although it may be acceptable for patient transport or community operations. Ambulance drivers may also have training in using the radio and knowing where medical supplies are stored in the ambulance.
Non-emergency attendant – This role has different levels of training across the world, but these staff are usually only required to perform patient transport duties (which can include stretcher or wheelchair cases), rather than acute care. Dependent on provider, they may be trained in first aid or extended skills such as use of an AED, oxygen therapy and other lifesaving or palliative skills. They may provide emergency cover when other units are not available, or when accompanied by a fully qualified technician or paramedic.
Emergency care assistant – Members of a frontline ambulance that drive the vehicles under both emergency and non-emergency conditions to incidents. Their role is to assist the clinician that they are working with, either a Technician or Paramedic, in their duties, whether that be drawing up drugs, setting up fluids (but not attaching), doing basic observations or performing 12 lead ECG assessments.
Emergency medical technician – technicians are usually able to perform a wide range of emergency care skills, such as defibrillation, spinal immobilization, bleeding control, splinting of suspected fractures, assisting the patient with certain medications, and oxygen therapy. Some countries split this term into levels (such as in the US, where there is EMT-Basic and EMT-Intermediate).
Registered nurse – In some systems, nurses are the primary providers of advanced-level care on ambulances, often in place of paramedics. This includes Estonia, the Netherlands, Sweden and Spain. Nurses may also work on ambulances for critical care transport.
Paramedic – This is a high level of medical training and usually involves key skills not permissible for technicians, such as cannulation (and with it the ability to administer a range of drugs such as morphine), tracheal intubation and other skills such as performing a cricothyrotomy. Dependent on jurisdiction, the title "paramedic" can be a protected title, and use of it without the relevant qualification may result in criminal prosecution.
Emergency care practitioner – This position is designed to bridge the link between ambulance care and the care of a general practitioner. ECPs are already qualified paramedics who have undergone further training, and are trained to prescribe medicines for longer-term care, such as antibiotics, as well as being trained in a range of additional diagnostic techniques.
Physician assistant – Physician Assistants are found predominately in English-speaking countries and may also be known as physician associates in some countries. PA's mirror the practice of a physician and are capable of providing the range of medical skills a physician provides. They generally work in collaboration with a physician, although in an ambulance environment this may not be possible. Instead, advanced directives or electronic communication is available to PA's to consult with physicians when required.
Physician – In some systems such as the SAMU in France, it is common for doctors to staff ambulances. On the other hand, this is rare in systems that rely heavily on paramedics or field nurses. In those cases, doctors may be present in specialist ambulance units – most notably the air ambulances. Alternatively, in some systems, such as Albuquerque, NM and Pittsburgh, PA, physicians are available to respond to serious cases via a fly car.
Military use
Military ambulances have historically included vehicles based on civilian designs and at times also included armored, but unarmed, vehicle ambulances based upon armoured personnel carriers (APCs). In the Second World War vehicles such as the Hanomag Sd Kfz 251 half-track were pressed into service as ad hoc ambulances, and in more recent times purpose-built AFVs such as the U.S. M1133 medical evacuation vehicle serve the exclusive purpose of armored medical vehicles. Civilian based designs may be painted in appropriate colors, depending on the operational requirements (i.e. camouflage for field use, white for United Nations peacekeeping, etc.). For example, the British Royal Army Medical Corps has a fleet of white ambulances, based on production trucks. Military helicopters have also served both as ad hoc and purpose-built air ambulances since they are extremely useful for MEDEVAC. In terms of equipment, military ambulances are barebones, often being nothing more than a box on wheels with racks to place manual stretchers, though for the operational conditions and level of care involved this is usually sufficient.
Since laws of war demand ambulances be marked with one of the Emblems of the Red Cross not to mount offensive weapons, military ambulances are often unarmed. It is a generally accepted practice in most countries to classify the personnel attached to military vehicles marked as ambulances as non-combatants; however, this does not always exempt medical personnel from coming under fireaccidental or deliberate. As a result, medics and other medical personnel attached to military ambulances are usually put through basic military training, on the assumption that they may have to use a weapon. The laws of war do allow non-combatant military personnel to carry individual weapons for protecting themselves and casualties. However, not all militaries exercise this right to their personnel.
The Israeli Defense Forces modified a number of its Merkava main battle tanks with ambulance features in order to allow rescue operations to take place under heavy fire in urban warfare. The modifications were made following a failed rescue attempt in which Palestinian gunmen killed two soldiers who were providing aid for a Palestinian woman in Rafah. Since M-113 armored personnel carriers and regular up-armored ambulances are not sufficiently protected against anti-tank weapons and improvised explosive devices, it was decided to use the heavily armored Merkava tank. Its rear door enables the evacuation of critically wounded soldiers. Israel did not remove the Merkava's weaponry, claiming that weapons were more effective protection than emblems since Palestinian militants would disregard any symbols of protection and fire at ambulances anyway. For use as ground ambulances and treatment & evacuation vehicles, the United States military currently employs the M113, the M577, the M1133 Stryker medical evacuation vehicle (MEV), and the RG-33 heavily armored ground ambulance (HAGA) as treatment and evacuation vehicles, with contracts to incorporate the newly designed M2A0 armored medical evacuation vehicle (AMEV), a variant of the M2 Bradley fighting vehicle (formerly known as the ATTV).
Some navies operate ocean-going hospital ships to lend medical assistance in high casualty situations such as wars or natural disasters. These hospital ships fulfill the criteria of an ambulance (transporting the sick or injured), although the capabilities of a hospital ship are more on par with a Mobile Army Surgical Hospital. In line with the laws of war, these ships can display a prominent Red Cross or Red Crescent to confer protection under the appropriate Geneva Convention. However, this designation has not always protected hospital ships from enemy fire.
Ambulette
Ambulettes provide patient transport service for non-emergency situations. Scheduling is a major factor in their effective use.
Reuse of retired ambulances
When an ambulance is retired, it may be donated or sold to another EMS provider. Alternately, it may be adapted into a storage and transport vehicle for crime scene identification equipment, a command post at community events, or support vehicle, such as a logistics unit. Others are refurbished and resold, or may just have their emergency equipment removed to be sold to private businesses or individuals, who then can use them as small recreational vehicles. They may also have a perfectly serviceable body or vehicle (or both) separated from the other and reused.
Toronto City Council operates a "Caravan of Hope" project to give retired Toronto ambulances a second life by donating them to the people of El Salvador. Since Ontario laws require ambulances to be retired after just four and a half years in service, the City of Toronto decommissions and auctions around 28 ambulances each year.
| Technology | Transport | null |
146351 | https://en.wikipedia.org/wiki/Hospital%20ship | Hospital ship | A hospital ship is a ship designated for primary function as a floating medical treatment facility or hospital. Most are operated by the military forces (mostly navies) of various countries, as they are intended to be used in or near war zones. In the 19th century, redundant warships were used as moored hospitals for seamen.
The Second Geneva Convention of 1949 prohibits military attacks on hospital ships that meet specified requirements, though belligerent forces have right of inspection and may take patients, but not staff, as prisoners of war.
History
Early examples
Hospital ships possibly existed in ancient times. The Athenian Navy had a ship named Therapia, and the Roman Navy had a ship named Aesculapius, their names indicating that they may have been hospital ships.
The earliest British hospital ship may have been the vessel Goodwill, which accompanied a Royal Navy squadron in the Mediterranean in 1608 and was used to house the sick sent aboard from other ships. However this experiment in medical care was short-lived, with Goodwill assigned to other tasks within a year and her complement of convalescents simply left behind at the nearest port. It was not until the mid-seventeenth century that any Royal Navy vessels were formally designated as hospital ships, and then only two throughout the fleet. These were either hired merchant ships or elderly sixth rates, with the internal bulkheads removed to create more room, and additional ports cut through the deck and hull to increase internal ventilation.
In addition to their sailing crew, these seventeenth century hospital ships were staffed by a surgeon and four surgeon's mates. The standard issue of medical supplies was bandages, soap, needles and bedpans. Patients were offered a bed or rug to rest upon, and given a clean pair of sheets. These early hospital ships were for the care of the sick rather than the wounded, with patients quartered according to their symptoms and infectious cases quarantined from the general population behind a sheet of canvas. The quality of food was very poor. In the 1690s, the surgeon aboard Siam complained that the meat was in an advanced state of putrefaction, the biscuits were weevil-ridden and bitter, and the bread was so hard that it stripped the skin off patients’ mouths.
Hospital ships were also used for the treatment of wounded soldiers fighting on land. An early example of this was during an English operation to evacuate English Tangier in 1683. An account of this evacuation was written by Samuel Pepys, an eyewitness. One of the main concerns was the evacuation of sick soldiers "and the many families and their effects to be brought off". The hospital ships Unity and Welcome sailed for England on 18 October 1683, with 114 invalid soldiers and 104 women and children, arriving at The Downs on 14 December 1683.
The number of medical personnel aboard Royal Navy hospital ships was slowly increased, with regulations issued in 1703 requiring that each vessel also carry six landsmen to act as surgical assistants, and four washerwomen. A 1705 amendment provided for a further five male nurses, and requisitions from the era suggest the number of sheets per patient was increased from one to two pairs. On 8 December 1798, unfit for service as a warship, was ordered to be converted to a hospital ship to hold wounded French and Spanish prisoners of war. According to Edward Hasted in 1798, two large hospital ships (also called lazarettos), (which were the surviving hulks of forty-four gun ships) were moored in Halstow Creek in Kent. The creek is an inlet from the River Medway and the River Thames. The crew of these vessels watched over ships coming to England, which were forced to stay in the creek under quarantine to protect the country from infectious diseases including the plague.
From 1821 to 1870, the Seamen's Hospital Society provided HMS Grampus, HMS Dreadnought and HMS Caledonia (later renamed Dreadnought) as successive hospital ships moored at Deptford in London. In 1866, HMS Hamadryad was moored in Cardiff as a seamen's hospital, replaced in 1905 by the Royal Hamadryad Seamen's Hospital. Other redundant warships were used as hospitals for convicts and prisoners of war.
Modern hospital ships
The Royal Navy institutionalised the use of hospital ships during the first half of the nineteenth century. Hospital ships were generally superior in their standard of service and sanitation to the medical provision available at the time for convalescent soldiers. The modern hospital ship began to emerge during the Crimean War in the 1850s. The only military hospital available to the British forces fighting on the Crimean Peninsula was at Scutari near the Bosphorus. During the Siege of Sevastopol almost 15,000 wounded troops were transported there from the port at Balaklava by a squadron of converted hospital ships.
The first ships to be equipped with genuine medical facilities were the steamships HMS Melbourne and HMS Mauritius, staffed by the Medical Staff Corps and providing services to the British expedition to China in 1860. The ships provided relatively spacious accommodation for the patients, and were equipped with an operating theatre. Another early hospital ship was in the 1860s, which aided the wounded soldiers of both sides during the American Civil War.
During the Russo-Turkish War (1877–78), the British Red Cross supplied a steel-hulled ship, equipped with modern surgery equipment including chloroform and other anaesthetics, and carbolic acid for antisepsis. Similar vessels accompanied the 1882 British invasion of Egypt and aided American personnel during the 1898 Spanish–American War.
During a smallpox outbreak in London in 1883, the Metropolitan Asylum Board (MAB) chartered and later purchased from the Admiralty two ships, and , and a paddle-steamer, , which were moored in the Thames at Long Reach, near Dartford, and remained in service until 1903.
Hospital ships were used by both sides in the Russo-Japanese War of 1904–1905.The sighting by the Japanese of the Russian hospital ship Orel, illuminated in accordance with regulations for hospital ships, led to the decisive naval Battle of Tsushima. Orel was retained as a prize of war by the Japanese after the battle.
World Wars
During World War I and World War II, hospital ships were first used on a massive scale. Many passenger liners were converted for use as hospital ships. and were two famous examples of ships serving in this capacity. By the end of the First World War, the British Royal Navy had 77 such ships in service. During the Gallipoli Campaign, hospital ships were used to evacuate wounded personnel to Egypt, Malta or England.
Canada operated hospital ships in both world wars. In World War I these included SS Letitia (I) and which was deliberately sunk by a German U-boat with great loss of life, despite the hospital ship's clearly marked status. In World War II, Canada operated the hospital ship and SS Letitia (II).
The first purpose-built hospital ship in the U.S. Navy was which was commissioned in 1921. During World War II both the United States Navy and Army operated hospital ships though with different purposes. Naval hospital ships were fully equipped hospitals designed to receive casualties direct from the battlefield and also supplied to provide logistical support to front line medical teams ashore. Army hospital ships were essentially hospital transports intended and equipped to evacuate patients from forward area Army hospitals to rear area hospitals or from those to the United States and were not equipped or staffed to handle large numbers of direct battle casualties. Three of the Navy hospital ships, , , and , were less elaborately equipped than other Navy hospital ships, medically staffed by Army medical personnel and similar in purpose to the Army model.
The last British royal yacht, the post World War II , was constructed in a way as to be convertible to a hospital ship in wartime. After her decommissioning, Peter Hennessy discovered that her actual role would have been as Queen Elizabeth II's refuge from nuclear weapons, hiding amidst the lochs of western Scotland.
A development of the Lun-class ekranoplan was planned for use as a mobile field hospital for rapid deployment to any ocean or coastal location at a speed of 297 knots (550 km/h, 341.8 mph). Work was 90% complete on this model, Spasatel, but Soviet military funding ceased and it was never completed.
Some hospital ships, such as and Esperanza del Mar, belong to civilian agencies, and do not belong to a navy. Mercy Ships is an international non-governmental charity (or NGO).
International law
Hospital ships were covered under the Hague Convention X of 1907. Articles of the Hague Convention X specified the provisions for a hospital ship:
Hospital-ships must be painted white. Military hospital ships must have a green band; ships operated by approved relief societies and similar must have a red band.
Ships must fly a red cross flag in addition to their national flag.
The ship should give medical assistance to wounded personnel of all nationalities.
The ship must not be used for any military purpose, or interfere with or hamper enemy combatant vessels.
Belligerents, as designated by the Hague Convention, can search any hospital ship to investigate violations of the above restrictions.
According to the San Remo Manual on International Law Applicable to Armed Conflicts at Sea, a hospital ship violating legal restrictions must be duly warned and given a reasonable time limit to comply. If a hospital ship persists in violating restrictions, a belligerent is legally entitled to capture it or take other means to enforce compliance. A non-complying hospital ship may only be fired on under the following conditions:
Diversion or capture is not feasible
No other method to exercise control is available
The violations are grave enough to allow the ship to be classified as a military objective
The damage and casualties will not be disproportionate to the military advantage.
In all other circumstances, attacking a hospital ship is a war crime.
Modern hospital ships display large Red Crosses or Red Crescents to signify their Geneva Convention protection under the laws of war. Even so, marked vessels have not been completely free from attack. Notable examples of hospital ships deliberately attacked during wartime are in 1915, the in 1941, and in 1943.
Current hospital ships
While any ship can be designated and marked as a hospital ship, many ships are permanently dedicated to that function.
Current military hospital ships
Current non-military hospital ships
Other shipborne hospitals
It is common for naval ships, especially large ships such as aircraft carriers and amphibious assault ships to have on-board hospitals. However, they are only one small part of the vessel's overall capability, and are used primarily for the ship's crew and its amphibious forces (and occasionally for relief missions). A warship with hospital facilities does not have the protected status of a hospital ship.
A primary example of the varied military-based hospital services available at sea is found aboard several types of US naval ships;
United States Navy;
– USS Gerald R. Ford, first in the class, has an on-board hospital that includes a full lab, pharmacy, operating room, 3-bed intensive care unit, 2-bed emergency room, and 41-bed hospital ward, staffed by 11 medical officers and 30 hospital corpsmen.
– Each carrier has a 53-bed hospital ward, a three-bed ICU, and acts as the hospital ship for the entire carrier strike group. In one year, the medical department of handled over 15,000 out-patient visits, drew almost 27,000 labs, filled almost 10,000 prescriptions, took about 2,300 x-rays and performed 65 surgical operations. There is not much variation among the ships of the class. The first ship, has 53 beds, plus 3 ICU beds, and the last ship, has 51 beds, plus 3 ICU beds.
(LHD) – These ships have 6 operating rooms, 14 ICU beds, 46 hospital beds, 4 battle dressing stations, medical imaging (i.e.:X-ray), a fully functional laboratory, and a blood bank. The ship can expand its medical complement to 600 beds, making it the second largest hospital at sea, second only to actual hospital ships.
amphibious assault ship (LHA) – This is the newest and largest class both in the USN and the world. However, the first two ships of the class, and , had the size of their medical facilities reduced, in favour of larger aviation facilities. The on-board hospitals of these first two vessels will have 2 operating rooms and 24 beds. It is unknown if this design change will affect the expanded capability for additional beds, nor what size the medical facilities of future ships of the class will be.
(LPD) – 24 hospital beds.
(LSD) – 11 hospital beds.
(LSD) – 8 hospital beds.
(EMS) - Will have four operating rooms and 124 medical beds, separated into acute care, acute isolation, ICU, and ICU isolation spaces.
More examples from various other national navies include;
Argentine Navy
– Icebreaker which can be deployed as a hospital ship.
Royal Australian Navy
– This class is based on the Juan Carlos I-class design, and has 2 operating rooms and a hospital ward.
People's Liberation Army Navy
Several armed s are fitted out as "ambulance transports".
Shichang – a multi-role training ship built in 1997. Deck space can accommodate modular medical units and can be used as a medical treatment facility, but the primary role is aviation training. The layout is very similar to RFA Argus (see below).
French Navy
– On board hospital is NATO Echelon level-3, with 69 hospital beds, 7 ICU beds, and an additional 50 beds if needed. The ship also has medical imaging capabilities, such as X-ray, CT-scan and ultrasound.
Italian Navy
aircraft carrier – Has an on-board hospital with 2 operating rooms, 1 intensive care unit, laboratory, pharmacy and a 32-bed hospital ward.
logistic ship – On-board hospital is NATO ROLE-level 2+, with operating room, intensive care unit and a laboratory.
Japan Maritime Self-Defense Force
– These ships have 2 operating rooms, 2 ICU beds, 35 hospital beds, 1 battle dressing station and several medical imaging (i.e.:X-ray) stations.
– These ships have 1 operating room, 1 ICU bed, 8 hospital beds.
- These ships have 1 operating room, 2 ICU beds, 6 hospital beds.
Spanish Navy
– Has a 40-bed hospital on board.
Royal Navy
Royal Fleet Auxiliary ship – This ship would be a hospital ship were it not for its armaments. However, it is instead designated as a 'Primary Casualty Receiving Ship' (PCRS). The vessel is classed as a NATO ROLE 3 Medical support vessel and is to be replaced in 2024
Royal Fleet Auxiliary Bay Class ships have a 14-bed medical facility which has the capability of being expanded in times of crisis as well as an operating theatre. The vessels are a classed as NATO Role 2 Medical support capable vessels.
German Navy
Berlin-class replenishment ship Berlin - Equipped with a container based version of the large modular hospital MERZ which stands for Marineeinsatzrettungszentrum (Englisch: Maritime Rescue Center) capable of holding 45 patients, plus 4 intensive care beds, clinical and microbiological laboratory and sterilisers.
Berlin-class replenishment ship Frankfurt am Main - Following a fire destroying the Frankfurt's MERZ, the Navy opted to equip the Frankfurt am Main with a new generation integrated MERZ (iMERZ), build into the hull of the ship. It's equipped with two operating rooms, medical imaging capabilities and a hospital ward. The German Navy plans to equip the Frankfurt's two sister ships with an iMERZ during routine maintenance.
| Technology | Naval transport | null |
146363 | https://en.wikipedia.org/wiki/Fire%20engine | Fire engine | A fire engine or fire truck (also spelled firetruck) is a vehicle, usually a specially-designed or modified truck, that functions as a firefighting apparatus. The primary purposes of a fire engine include transporting firefighters and water to an incident as well as carrying equipment for firefighting operations in a fire drill. Some fire engines have specialized functions, such as wildfire suppression and aircraft rescue and firefighting, and may also carry equipment for technical rescue.
Many fire engines are based on a commercial vehicle chassis that is further upgraded and customized for firefighting requirements. They are generally considered emergency vehicles authorized to be equipped with emergency lights and sirens, as well as communication equipment such as two-way radios and mobile computer technology.
The terms fire engine and fire truck are often used interchangeably to a broad range of vehicles involved in firefighting; however, in some fire departments they refer to separate and specific types of vehicle.
Design and construction
The design and construction of fire engines focuses greatly on the use of both active and passive warnings. Passive visual warnings involve the use of high contrast patterns to increase the noticeability of the vehicle. These types of warnings are often seen on older vehicles and those in developing countries. More modern designs make use of retroreflectors to reflect light from other vehicles. Vehicles will also often have these reflectors arranged in a chevron pattern along with the words fire or rescue. European countries commonly use a pattern known as Battenburg markings. Along with the passive warnings, are active visual warnings which are usually in the form of flashing colored lights (also known as "beacons" or "lightbars"). These flash to attract the attention of other drivers as the fire truck approaches, or to provide warning to drivers approaching a parked fire truck in a dangerous position on the road. While the fire truck is headed towards the scene, the lights are always accompanied by loud audible warnings such as sirens and air horns. Some fire engines in the United States are lime yellow rather than red due to safety and ergonomics reasons. A 2009 study by the U.S. Fire Administration concluded that fluorescent colors, including yellow-green and orange, are easiest to spot in daylight.
In some regions, a fire engine may be used to transport first responder firefighters, paramedics or EMTs to medical emergencies due to their proximity to the incident.
Types
Conventional fire engine
The standard fire engine transports firefighters to the scene, carries equipment needed by the firefighters for most firefighting scenarios, and may provide a limited supply of water with which to fight the fire. The tools carried on the fire engine will vary greatly based on many factors including the size of the department and the usual situations the firefighters handle. For example, departments located near large bodies of water or rivers are likely to have some sort of water rescue equipment. Standard tools found on nearly all fire engines include ladders, hydraulic rescue tools (often referred to as the jaws of life), floodlights, fire hose, fire extinguishers, self-contained breathing apparatus, and thermal imaging cameras.
The exact layout of what is carried on an engine is decided by the needs of the department. For example, fire departments located in metropolitan areas will carry equipment to mitigate hazardous materials and effect technical rescues, while departments that operate in the wildland-urban interface will need the gear to deal with brush fires.
Some fire engines have a fixed deluge gun, also known as a master stream, which directs a heavy stream of water to wherever the operator points it. An additional feature of engines are their preconnected hose lines, commonly referred to as preconnects. The preconnects are attached to the engine's onboard water supply and allow firefighters to quickly mount an aggressive attack on the fire as soon as they arrive on scene. When the onboard water supply runs out, the engine is connected to more permanent sources such as fire hydrants or water tenders and can also use natural sources such as rivers or reservoirs by drafting water.
Aerial apparatus
An aerial apparatus is a fire truck mounted with an extendable boom that enables firefighters to reach high locations. They can provide a high vantage point for spraying water and creating ventilation, an access route for firefighters and an escape route for firefighters and people they have rescued. In North America, aerial apparatuses are used for fire suppression, whereas in Europe, they are used more for rescue.
Turntable ladder
A turntable ladder (TL) is an aerial apparatus with a large ladder mounted on a pivot which resembles a turntable, giving it its name. The key functions of a turntable ladder are allowing access or egress of firefighters and fire victims at height, providing a high-level water point for firefighting (elevated master stream), and providing a platform from which tasks such as ventilation or overhaul can be executed.
To increase its length and reach, the ladder is often telescoping. Modern telescopic ladders may be hydraulic or pneumatic. These mechanical features allow the use of ladders which are longer, sturdier, and more stable. They may also have pre-attached hoses or other equipment.
The pivot can be mounted at the rear of the chassis or in the middle, just behind the cab. The latter is sometimes called a "mid-ship" arrangement, and it allows a lower travel height for the truck.
While the traditional characteristic of a TL was a lack of water pumping or storage, many modern TLs have a water pumping function built in (and some have their own on-board supply reservoir). Some may have piping along the ladder to supply water to firefighters at the top of the ladder, and some of these may also have a monitor installed at the top. Other appliances may simply have a track-way to securely hold a manually-run hose reel.
In the United States, turntable ladders with additional functions such as an onboard pump, a water tank, fire hose, aerial ladder and multiple ground ladders, are known as quad or quint engines, indicating the number of functions they perform.
The highest TL in the world is the Magirus M68L, with a range of .
Tiller truck
In the United States, a tiller truck, also known as a tractor-drawn aerial (TDA), tiller ladder, or hook-and-ladder truck, is a specialized turntable ladder mounted on a semi-trailer truck. Unlike a commercial semi, the trailer and tractor are permanently combined and special tools are required to separate them. It has two drivers, with separate steering wheels for front and rear wheels.
One of the main features of the tiller-truck is its enhanced maneuverability. The independent steering of the front and back wheels allow the tiller to make much sharper turns, which is particularly helpful on narrow streets and in apartment complexes with maze-like roads. An additional feature of the tiller-truck is that its overall length, over for most models, allows for additional storage of tools and equipment. The extreme length gives compartment capacities that range between in the trailer with an additional in the cab.
Some departments elect to use tiller-quints, which are tiller trucks that have the added feature of being fitted with an on-board water tank. These are particularly useful for smaller departments that do not have enough personnel to staff both an engine company and a truck company.
Platform truck
A platform truck carries an aerial work platform, also known as a basket or bucket, on the end of a ladder or boom. These platforms can provide a secure place from which a firefighter can operate. Many platforms also allow for rescues to be performed and are outfitted with tie down clips and rappelling arms.
Some booms are capable of articulating, allowing the arm to bend in one or more places. This allows the platform truck to go "up and over" an obstacle, and is an advantage over the traditional platform ladder, which can only extend in a straight line.
Wildland fire engine
A wildland fire engine is a specialized fire engine that can negotiate difficult terrain for wildfire suppression. A wildland fire engine is smaller than standard fire engines and has a higher ground clearance. They may also respond to emergencies in rough terrain where other vehicles cannot respond. Many wildland engines feature four-wheel drive capability to improve hill climbing and rough terrain capability. Some wildland apparatus can pump water while driving (compared to some traditional engines which must be stationary to pump water), allowing "mobile attacks" on vegetation fires to minimize the rate of spread.
Fire departments that serve areas along the wildland–urban interface have to be able to tackle traditional urban fires as well as wildland fires. Departments in these areas often use a wildland-urban interface engine, which combine features of a standard fire engine with that of a wildland fire engine.
Water tender
A water tender is a specialist fire appliance with the primary purpose of transporting large amounts of water to the fire area to make it available for extinguishing operations. These are especially useful in rural areas where fire hydrants are not readily available and natural water resources are insufficient or difficult to exploit.
Most tankers have an on-board pumping system. This pump is often not of sufficient power to fight fires (as it is designed to be attached to a fire engine), but is more often used to draw water into the tender from hydrants or other water sources. Many tankers are equipped with fast-drain valves on the sides and back of the truck. This allows firefighters to empty thousands of gallons of water into a portable water tank in just a few seconds.
Most water tenders are designed to carry loads of .
Airport crash tender
An airport crash tender is a specialized fire engine designed for use at aerodromes in aircraft accidents. Some of the features that make the airport crash tender unique are its ability to move on rough terrain outside the runway and airport area, large water capacity as well as a foam tank, a high-capacity pump, and water/foam monitors. Newer airport crash tenders also incorporate twin agent nozzles/injection systems that add dry chemical fire retardant (such as Purple-K) to create a stream of firefighting foam which is able to stop the fire faster. Some also have gaseous fire suppression tanks for electrical fires. These features give the airport crash tenders a capability to reach an airplane rapidly, and rapidly extinguish large fires with jet fuel involved.
Other vehicles
Other vehicles that are used by fire departments but may not be directly involved in firefighting may include
Fire car
Fire investigation unit
Fire police unit
Hazardous materials apparatus
Light and air unit
Marine rescue unit
Mobile communications vehicle
Operational support unit
History
An early device used to squirt water onto a fire was known as a squirt or fire syringe. Hand squirts and hand pumps are noted before Ctesibius of Alexandria invented the first fire pump around the 2nd century B.C., and an example of a force-pump possibly used for a fire-engine is mentioned by Heron of Alexandria.
In 1650, Hans Hautsch built a fire engine with a compressed air vessel. On each side 14 men worked a piston rod back and forth in a horizontal direction. The air vessel, a type of pressure tank, issued an even stream despite the backward motion of the piston. This was made possible by a rotating pipe mounted on the hose which allowed the jet to reach heights up to . Caspar Schott observed Hautsch's fire engine in 1655 and wrote an account of it in his Magia Universalis.
Colonial laws in America required each house to have a bucket of water on the front stoop in preparation for fires at night. These buckets were intended for use by the initial bucket brigade that would supply the water at fires. Philadelphia obtained a hand-pumped fire engine in 1719, years after Boston's 1654 model appeared there, made by Joseph Jenckes Sr., but before New York's two engines arrived from London.
By 1730, Richard Newsham, in London, had made successful fire engines. He also invented those first used in New York City in 1731 where the amount of manpower and skill necessary for firefighting prompted Benjamin Franklin to found an organized fire company in 1737. Thomas Lote built the first fire engine made in America in 1743. These earliest engines are called hand tubs because they are manually (hand) powered and the water was supplied by a bucket brigade dumping it into a tub (cistern) where the pump had a permanent intake pipe.
An important advancement around 1822 was the invention of an engine which could draft water from a water source. This rendered the bucket brigade obsolete. In 1822, a Philadelphia-based manufacturing company called Sellers and Pennock made a model called "The Hydraulion". It is said to be the first suction engine. Some models had the hard, suction hose fixed to the intake and curled up over the apparatus known as a squirrel tail engine.
The earliest engines were small and were either carried by four men, or mounted on skids and dragged to a fire. As the engines grew larger they became horse-drawn and later self-propelled by steam engines.
Until the mid-19th century, most fire engines were maneuvered by men, but the introduction of horse-drawn fire engines considerably improved the response time to incidents. The first self-propelled steam pumper fire engine was built in New York in 1841. Unfortunately for the manufacturers, some firefighters sabotaged the device and its use of the first engine was discontinued. However, the need and the utility of power equipment ensured the success of the steam pumper well into the twentieth century. Many cities and towns around the world bought the steam fire engines.
Motorised fire engines date back to January 1897, when the Prefect of Police in Paris applied for funds to purchase "a machine worked by petroleum for the traction of a fire-engine, ladders, and so forth and for the conveyance of the necessary staff of pompiers". With great prescience the report states "If the experiment prove successful, as is anticipated, horses will eventually be entirely replaced by automobiles". This was, indeed, the case and motorised fire engines became commonplace by the early 20th century. By 1905, the idea of combining gas engine motor trucks into fire engines was attracting great attention; according to a Popular Mechanics article in that year, such trucks were rapidly gaining popularity in England. That same year, the Knox Automobile Company of Springfield, Massachusetts, began selling what some have described as the world's first modern fire engine. A year later, the city of Springfield, Illinois, had filled their fire department with Knox engines. Another early motorized fire engine was developed by Peter Pirsch and Sons of Kenosha, Wisconsin.
For many years firefighters sat on the sides of the fire engines, or even stood on the rear of the vehicles, exposed to the elements. This arrangement was uncomfortable and dangerous (some firefighters were thrown to their deaths when their fire engines made sharp turns on the road), and today nearly all fire engines have fully enclosed seating areas for their crews.
Hook and ladder
The "hook and ladder" was an early type of fire units known since late 1700s. It was a horse-drawn carriage which brought ladders and hooks to the fire place. Ladders were used for access to upper floors and the roof. "Hooks" were pike poles used for pulling down and apart the burning construction.
Early pumpers
Early pumpers used cisterns as a source of water. Water was later put into wooden pipes under the streets and a "fire plug" was pulled out of the top of the pipe when a suction hose was to be inserted. Later systems incorporated pressurized fire hydrants, where the pressure was increased when a fire alarm was sounded. This was found to be harmful to the system and unreliable. Today's valved hydrant systems are kept under pressure at all times, although additional pressure may be added when needed. Pressurized hydrants eliminate much of the work in obtaining water for pumping through the engine and into the attack hoses. Many rural fire engines still rely upon cisterns or other sources for drafting water into the pumps. Steam pumper came in to use in the 1850s.
Early aerials
In the late 19th century, means of reaching tall structures were devised. At first, manually extendable ladders were used; as these grew in length (and weight), they were put onto two large wheels. When carried by fire engines these wheeled escape ladders had the wheels suspended behind the rear of the vehicle, making them a distinctive sight. Before long, turntable ladders—which were even longer, mechanically extendable, and installed directly onto fire trucks—made their appearances.
After World War II, turntable ladders were supplemented by the aerial work platform (sometimes called "cherry picker"), a platform or bucket attached onto a mechanically bending arm (or "snorkel") installed onto a fire truck. While these could not reach the height of similar turntable ladders, the platforms could extend into previously unreachable "dead corners" of a burning building.
| Technology | Specific-purpose transportation | null |
146384 | https://en.wikipedia.org/wiki/Defibrillation | Defibrillation | Defibrillation is a treatment for life-threatening cardiac arrhythmias, specifically ventricular fibrillation (V-Fib) and non-perfusing ventricular tachycardia (V-Tach). A defibrillator delivers a dose of electric current (often called a counter-shock) to the heart. Although not fully understood, this process depolarizes a large amount of the heart muscle, ending the arrhythmia. Subsequently, the body's natural pacemaker in the sinoatrial node of the heart is able to re-establish normal sinus rhythm. A heart which is in asystole (flatline) cannot be restarted by a defibrillator; it would be treated only by cardiopulmonary resuscitation (CPR) and medication, and then by cardioversion or defibrillation if it converts into a shockable rhythm.
In contrast to defibrillation, synchronized electrical cardioversion is an electrical shock delivered in synchrony to the cardiac cycle. Although the person may still be critically ill, cardioversion normally aims to end poorly perfusing cardiac arrhythmias, such as supraventricular tachycardia.
Defibrillators can be external, transvenous, or implanted (implantable cardioverter-defibrillator), depending on the type of device used or needed. Some external units, known as automated external defibrillators (AEDs), automate the diagnosis of treatable rhythms, meaning that lay responders or bystanders are able to use them successfully with little or no training.
Use of defibrillators
Indications
Defibrillation is often an important step in cardiopulmonary resuscitation (CPR). CPR is an algorithm-based intervention aimed to restore cardiac and pulmonary function. Defibrillation is indicated only in certain types of cardiac dysrhythmias, specifically ventricular fibrillation (VF) and pulseless ventricular tachycardia. If the heart has completely stopped, as in asystole or pulseless electrical activity (PEA), defibrillation is not indicated. Defibrillation is also not indicated if the patient is conscious or has a pulse. Improperly given electrical shocks can cause dangerous dysrhythmias, such as ventricular fibrillation.
Application method
A defibrillation device that is often available outside of medical centers is the automated external defibrillator (AED), a portable machine that can be used with no previous training. That is possible because the machine produces pre-recorded voice instructions that guide the user. The device automatically checks the patient's condition and applies the correct electric shocks. There also exist written instructions that explain the procedure step-by-step.
Outcomes
Survival rates for out-of-hospital cardiac arrests in North America are poor, often less than 10%. Outcome for in-hospital cardiac arrests are higher at 20%. Within the group of people presenting with cardiac arrest, the specific cardiac rhythm can significantly impact survival rates. Compared to people presenting with a non-shockable rhythm (such as asystole or PEA), people with a shockable rhythm (such as VF or pulseless ventricular tachycardia) have improved survival rates, ranging between 21 and 50%.
Types
Manual models
Manual external defibrillators require the expertise of a healthcare professional. They are used in conjunction with an electrocardiogram, which can be separate or built-in. A healthcare provider first diagnoses the cardiac rhythm and then manually determine the voltage and timing for the electrical shock. These units are primarily found in hospitals and on some ambulances. For instance, every NHS ambulance in the United Kingdom is equipped with a manual defibrillator for use by the attending paramedics and technicians. In the United States, many advanced EMTs and all paramedics are trained to recognize lethal arrhythmias and deliver appropriate electrical therapy with a manual defibrillator when appropriate.
An internal defibrillator is often used to defibrillate the heart during or after cardiac surgery such as a heart bypass. The electrodes consist of round metal plates that come in direct contact with the myocardium. Manual internal defibrillators deliver the shock through paddles placed directly on the heart. They are mostly used in the operating room and, in rare circumstances, in the emergency room during an open heart procedure.
Automated external defibrillators
Automated external defibrillators (AEDs) are designed for use by untrained or briefly trained laypersons. AEDs contain technology for analysis of heart rhythms. As a result, it does not require a trained health provider to determine whether or not a rhythm is shockable. By making these units publicly available, AEDs have improved outcomes for sudden out-of-hospital cardiac arrests.
Trained health professionals have more limited use for AEDs than manual external defibrillators. Recent studies show that AEDs does not improve outcome in patients with in-hospital cardiac arrests. AEDs have set voltages and does not allow the operator to vary voltage according to need. AEDs may also delay delivery of effective CPR. For diagnosis of rhythm, AEDs often require the stopping of chest compressions and rescue breathing. For these reasons, certain bodies, such as the European Resuscitation Council, recommend using manual external defibrillators over AEDs if manual external defibrillators are readily available.
As early defibrillation can significantly improve VF outcomes, AEDs have become publicly available in many easily accessible areas. AEDs have been incorporated into the algorithm for basic life support (BLS). Many first responders, such as firefighters, police officers, and security guards, are equipped with them.
AEDs can be fully automatic or semi-automatic. A semi-automatic AED automatically diagnoses heart rhythms and determines if a shock is necessary. If a shock is advised, the user must then push a button to administer the shock. A fully automated AED automatically diagnoses the heart rhythm and advises the user to stand back while the shock is automatically given. Some types of AEDs come with advanced features, such as a manual override or an ECG display.
Cardioverter-defibrillators
Implantable cardioverter-defibrillators, also known as automatic internal cardiac defibrillator (AICD), are implants similar to pacemakers (and many can also perform the pacemaking function). They constantly monitor the patient's heart rhythm, and automatically administer shocks for various life-threatening arrhythmias, according to the device's programming. Many modern devices can distinguish between ventricular fibrillation, ventricular tachycardia, and more benign arrhythmias like supraventricular tachycardia and atrial fibrillation. Some devices may attempt overdrive pacing prior to synchronised cardioversion. When the life-threatening arrhythmia is ventricular fibrillation, the device is programmed to proceed immediately to an unsynchronized shock.
There are cases where the patient's ICD may fire constantly or inappropriately. This is considered a medical emergency, as it depletes the device's battery life, causes significant discomfort and anxiety to the patient, and in some cases may actually trigger life-threatening arrhythmias. Some emergency medical services personnel are now equipped with a ring magnet to place over the device, which effectively disables the shock function of the device while still allowing the pacemaker to function (if the device is so equipped). If the device is shocking frequently, but appropriately, EMS personnel may administer sedation.
A wearable cardioverter defibrillator is a portable external defibrillator that can be worn by at-risk patients. The unit monitors the patient 24 hours a day and can automatically deliver a biphasic shock if VF or VT is detected. This device is mainly indicated in patients who are not immediate candidates for ICDs.
Interface
The connection between the defibrillator and the patient consists of a pair of electrodes, each provided with electrically conductive gel in order to ensure a good connection and to minimize electrical resistance, also called chest impedance (despite the DC discharge) which would burn the patient. Gel may be either wet (similar in consistency to surgical lubricant) or solid (similar to gummi candy). Solid-gel is more convenient, because there is no need to clean the used gel off the person's skin after defibrillation. However, the use of solid-gel presents a higher risk of burns during defibrillation, since wet-gel electrodes more evenly conduct electricity into the body. Paddle electrodes, which were the first type developed, come without gel, and must have the gel applied in a separate step. Self-adhesive electrodes come prefitted with gel. There is a general division of opinion over which type of electrode is superior in hospital settings; the American Heart Association favors neither, and all modern manual defibrillators used in hospitals allow for swift switching between self-adhesive pads and traditional paddles. Each type of electrode has its merits and demerits.
Paddle electrodes
The most well-known type of electrode (widely depicted in films and television) is the traditional metal "hard" paddle with an insulated (usually plastic) handle. This type must be held in place on the patient's skin with approximately 25 lbs (11.3 kg) of force while a shock or a series of shocks is delivered. Paddles offer a few advantages over self-adhesive pads. Many hospitals in the United States continue the use of paddles, with disposable gel pads attached in most cases, due to the inherent speed with which these electrodes can be placed and used. This is critical during cardiac arrest, as each second of nonperfusion means tissue loss. Modern paddles allow for monitoring (electrocardiography), though in hospital situations, separate monitoring leads are often already in place.
Paddles are reusable, being cleaned after use and stored for the next patient. Gel is therefore not preapplied, and must be added before these paddles are used on the patient. Paddles are generally only found on manual external units.
Self-adhesive electrodes
Newer types of resuscitation electrodes are designed as an adhesive pad, which includes either solid or wet gel. These are peeled off their backing and applied to the patient's chest when deemed necessary, much the same as any other sticker. The electrodes are then connected to a defibrillator, much as the paddles would be. If defibrillation is required, the machine is charged, and the shock is delivered, without any need to apply any additional gel or to retrieve and place any paddles. Most adhesive electrodes are designed to be used not only for defibrillation, but also for transcutaneous pacing and synchronized electrical cardioversion. These adhesive pads are found on most automated and semi-automated units and are replacing paddles entirely in non-hospital settings. In hospital, for cases where cardiac arrest is likely to occur (but has not yet), self-adhesive pads may be placed prophylactically.
Pads also offer an advantage to the untrained user, and to medics working in the sub-optimal conditions of the field. Pads do not require extra leads to be attached for monitoring, and they do not require any force to be applied as the shock is delivered. Thus, adhesive electrodes minimize the risk of the operator coming into physical (and thus electrical) contact with the patient as the shock is delivered by allowing the operator to be up to several feet away. (The risk of electrical shock to others remains unchanged, as does that of shock due to operator misuse.) Self-adhesive electrodes are single-use only. They may be used for multiple shocks in a single course of treatment, but are replaced if (or in case) the patient recovers then reenters cardiac arrest.
Special pads are used for children under the age of 8 or those under 55 lbs. (22 kg).
Placement
Resuscitation electrodes are placed according to one of two schemes. The anterior-posterior scheme is the preferred scheme for long-term electrode placement. One electrode is placed over the left precordium (the lower part of the chest, in front of the heart). The other electrode is placed on the back, behind the heart in the region between the scapula. This placement is preferred because it is best for non-invasive pacing.
The anterior-apex scheme (anterior-lateral position) can be used when the anterior-posterior scheme is inconvenient or unnecessary. In this scheme, the anterior electrode is placed on the right, below the clavicle. The apex electrode is applied to the left side of the patient, just below and to the left of the pectoral muscle. This scheme works well for defibrillation and cardioversion, as well as for monitoring an ECG.
Researchers have created a software modeling system capable of mapping an individual's chest and determining the best position for an external or internal cardiac defibrillator.
Mechanism
Defibrillation halts chaotic cardiac activity by forcibly depolarizing heart cells, disrupting re-entrant circuits, and allowing for the heart's natural pacemaker to take over.
Cardiac cells require a strong electrical stimulus to raise their transmembrane potential to the activation threshold. Only a small amount of electrical current enters the cell due to high membrane impedance.The intracellular voltage of the cell remains uniform, while the extracellular voltage rapidly increases or decreases depending on proximity to the electrodes.This creates a voltage gradient that alters the transmembrane potential of cells, potentially resetting irregular electrical activity to restore normal cardiac rhythm.
Irregular rhythms often result from re-entrant circuits, where electrical impulses circle within the heart tissue due to areas of slow conduction or unidirectional block. The widespread depolarization from the shock interrupts these circuits, stopping the erratic propagation of electrical signals.
After the cells depolarize, they enter a refractory period, during which they cannot be re-excited. This allows the heart's natural pacemaker, the sinoatrial node, to resume control of the rhythm. During this period, ion pumps actively restore the normal distribution of ions, re-establishing the resting membrane potential.
History
Defibrillators were first demonstrated in 1899 by Jean-Louis Prévost and Frédéric Batelli, two physiologists from the University of Geneva, Switzerland. They discovered that small electrical shocks could induce ventricular fibrillation in dogs, and that larger charges would reverse the condition.
In 1933, Dr. Albert Hyman, heart specialist at the Beth Davis Hospital of New York City and C. Henry Hyman, an electrical engineer, looking for an alternative to injecting powerful drugs directly into the heart, came up with an invention that used an electrical shock in place of drug injection. This invention was called the Hyman Otor where a hollow needle is used to pass an insulated wire to the heart area to deliver the electrical shock. The hollow steel needle acted as one end of the circuit and the tip of the insulated wire the other end. Whether the Hyman Otor was a success is unknown.
The external defibrillator, as it is known today, was invented by electrical engineer William Kouwenhoven in 1930. Kouwenhoven studied the relationship between electric shocks and their effects on the human heart when he was a student at Johns Hopkins University School of Engineering. His studies helped him invent a device to externally jump start the heart. He invented the defibrillator and tested it on a dog, like Prévost and Batelli. The first use on a human was in 1947 by Claude Beck, professor of surgery at Case Western Reserve University.
Beck's theory was that ventricular fibrillation often occurred in hearts that were fundamentally healthy, in his terms "Hearts that are too good to die", and that there must be a way of saving them. Beck first used the technique successfully on a 14-year-old boy who was having his breastbone separated from his ribs because of a congenital growth disorder, causing breathing problems. The boy's chest was surgically opened, and manual cardiac massage was undertaken for 45 minutes until the arrival of the defibrillator. Beck used internal paddles on either side of the heart, along with procainamide, an antiarrhythmic drug, and achieved return of a perfusing cardiac rhythm.
These early defibrillators used the alternating current from a power socket, transformed from the 110–240 volts available in the line, up to between 300 and 1000 volts, to the exposed heart by way of "paddle" type electrodes. The technique was often ineffective in reverting VF while morphological studies showed damage to the cells of the heart muscle post-mortem. The nature of the AC machine with a large transformer also made these units very hard to transport, and they tended to be large units on wheels.
Closed-chest method
Until the early 1950s, defibrillation of the heart was possible only when the chest cavity was open during surgery. The technique used an alternating voltage from a 300 or greater volt source derived from standard AC power, delivered to the sides of the exposed heart by "paddle" electrodes where each electrode was a flat or slightly concave metal plate of about 40 mm diameter. The closed-chest defibrillator device which applied an alternating voltage of greater than 1000 volts, conducted by means of externally applied electrodes through the chest cage to the heart, was pioneered by Dr V. Eskin with assistance by A. Klimov in Frunze, USSR (today known as Bishkek, Kyrgyzstan) in the mid-1950s. The duration of AC shocks was typically in the range of 100–150 milliseconds.
Direct current method
Early successful experiments of successful defibrillation by the discharge of a capacitor performed on animals were reported by N. L. Gurvich and G. S. Yunyev in 1939. In 1947 their works were reported in western medical journals. Serial production of Gurvich's pulse defibrillator started in 1952 at the electromechanical plant of the institute, and was designated model ИД-1-ВЭИ (Импульсный Дефибриллятор 1, Всесоюзный Электротехнический Институт, or in English, Pulse Defibrillator 1, All-Union Electrotechnical Institute). It is described in detail in Gurvich's 1957 book, Heart Fibrillation and Defibrillation.
The first Czechoslovak "universal defibrillator Prema" was manufactured in 1957 by the company Prema, designed by Dr. Bohumil Peleška. In 1958 his device was awarded Grand Prix at Expo 58.
In 1958, US senator Hubert H. Humphrey visited Nikita Khrushchev and among other things he visited the Moscow Institute of Reanimatology, where, among others, he met with Gurvich. Humphrey immediately recognized importance of reanimation research and after that a number of American doctors visited Gurvich. At the same time, Humphrey worked on establishing a federal program in the National Institute of Health in physiology and medicine, telling Congress: "Let's compete with U.S.S.R. in research on reversibility of death".
In 1959 Bernard Lown commenced research in his animal laboratory in collaboration with engineer Barouh Berkovits into a technique which involved charging of a bank of capacitors to approximately 1000 volts with an energy content of 100–200 joules then delivering the charge through an inductance such as to produce a heavily damped sinusoidal wave of finite duration (~5 milliseconds) to the heart by way of paddle electrodes. This team further developed an understanding of the optimal timing of shock delivery in the cardiac cycle, enabling the application of the device to arrhythmias such as atrial fibrillation, atrial flutter, and supraventricular tachycardias in the technique known as "cardioversion".
The Lown-Berkovits waveform, as it was known, was the standard for defibrillation until the late 1980s. Earlier in the 1980s, the "MU lab" at the University of Missouri had pioneered numerous studies introducing a new waveform called a biphasic truncated waveform (BTE). In this waveform an exponentially decaying DC voltage is reversed in polarity about halfway through the shock time, then continues to decay for some time after which the voltage is cut off, or truncated. The studies showed that the biphasic truncated waveform could be more efficacious while requiring the delivery of lower levels of energy to produce defibrillation. An added benefit was a significant reduction in weight of the machine. The BTE waveform, combined with automatic measurement of transthoracic impedance, is the basis for modern defibrillators.
Portable units
A major breakthrough was the introduction of portable defibrillators used out of the hospital. Already Peleška's Prema defibrillator was designed to be more portable than original Gurvich's model. In Soviet Union, a portable version of Gurvich's defibrillator, model ДПА-3 (DPA-3), was reported in 1959. In the west this was pioneered in the early 1960s by Prof. Frank Pantridge in Belfast. Today portable defibrillators are among the many very important tools carried by ambulances. They are the only proven way to resuscitate a person who has had a cardiac arrest unwitnessed by Emergency Medical Services (EMS) who is still in persistent ventricular fibrillation or ventricular tachycardia at the arrival of pre-hospital providers.
Gradual improvements in the design of defibrillators, partly based on the work developing implanted versions (see below), have led to the availability of Automated External Defibrillators. These devices can analyse the heart rhythm by themselves, diagnose the shockable rhythms, and charge to treat. This means that no clinical skill is required in their use, allowing lay people to respond to emergencies effectively.
Waveform change
Until the mid 1990s, external defibrillators delivered a Lown type waveform (see Bernard Lown), a heavily damped sinusoidal impulse having a mainly uniphasic characteristic. Biphasic defibrillation alternates the direction of the pulses, completing one cycle in approximately 12 milliseconds. Biphasic defibrillation was originally developed and used for implantable cardioverter-defibrillators. When applied to external defibrillators, biphasic defibrillation significantly decreases the energy level necessary for successful defibrillation, decreasing the risk of burns and myocardial damage.
Ventricular fibrillation (VF) could be returned to sinus rhythm in 60% of cardiac arrest patients treated with a single shock from a monophasic defibrillator. Most biphasic defibrillators have a first shock success rate of greater than 90%.
Implantable devices
A further development in defibrillation came with the invention of the implantable device, known as an implantable cardioverter-defibrillator (or ICD). This was pioneered at Sinai Hospital in Baltimore by a team that included Stephen Heilman, Alois Langer, Jack Lattuca, Morton Mower, Michel Mirowski, and Mir Imran, with the help of industrial collaborator Intec Systems of Pittsburgh. Mirowski teamed up with Mower and Staewen, and together they commenced their research in 1969. However, it was 11 years before they treated their first patient. Similar developmental work was carried out by Schuder and colleagues at the University of Missouri.
The work was commenced, despite doubts amongst leading experts in the field of arrhythmias and sudden death. There was doubt that their ideas would ever become a clinical reality. In 1962 Bernard Lown introduced the external DC defibrillator. This device applied a direct current from a discharging capacitor through the chest wall into the heart to stop heart fibrillation. In 1972, Lown stated in the journal Circulation – "The very rare patient who has frequent bouts of ventricular fibrillation is best treated in a coronary care unit and is better served by an effective antiarrhythmic program or surgical correction of inadequate coronary blood flow or ventricular malfunction. In fact, the implanted defibrillator system represents an imperfect solution in search of a plausible and practical application."
The problems to be overcome were the design of a system which would allow detection of ventricular fibrillation or ventricular tachycardia. Despite the lack of financial backing and grants, they persisted and the first device was implanted in February 1980 at Johns Hopkins Hospital by Dr. Levi Watkins Jr. assisted by Vivien Thomas. Modern ICDs do not require a thoracotomy and possess pacing, cardioversion, and defibrillation capabilities.
The invention of implantable units is invaluable to some people with regular heart problems, although they are generally only given to those people who have already had a cardiac episode.
People can live long normal lives with the devices. Many patients have multiple implants. A patient in Houston, Texas had an implant at the age of 18 in 1994 by the recent Dr. Antonio Pacifico. He was awarded "Youngest Patient with Defibrillator" in 1996. Today these devices are implanted into small babies shortly after birth.
Society and culture
As devices that can quickly produce dramatic improvements in patient health, defibrillators are often depicted in movies, television, video games and other fictional media. Their function, however, is often exaggerated with the defibrillator inducing a sudden, violent jerk or convulsion by the patient. The pad placement is also shown wrong, along with sudden rising of patient to large height when shock is given. In reality, while the muscles may contract, such dramatic patient presentation is rare. Similarly, medical providers are often depicted defibrillating patients with a "flat-line" ECG rhythm (also known as asystole). This is not normal medical practice, as the heart cannot be restarted by the defibrillator itself. Only the cardiac arrest rhythms ventricular fibrillation and pulseless ventricular tachycardia are normally defibrillated. The purpose of defibrillation is to depolarize the entire heart all at once so that it is synchronized, effectively inducing temporary asystole, in the hope that in the absence of the previous abnormal electrical activity, the heart will spontaneously resume beating normally. Someone who is already in asystole cannot be helped by electrical means, and usually needs urgent CPR and intravenous medication (and even these are rarely successful in cases of asystole). A useful analogy to remember is to think of defibrillators as power-cycling, rather than jump-starting, the heart. There are also several heart rhythms that can be "shocked" when the patient is not in cardiac arrest, such as supraventricular tachycardia and ventricular tachycardia that produces a pulse; this more-complicated procedure is known as cardioversion, not defibrillation.
In Australia up until the 1990s it was relatively rare for ambulances to carry defibrillators. This changed in 1990 after Australian media mogul Kerry Packer had a cardiac arrest due to a heart attack and, purely by chance, the ambulance that responded to the call carried a defibrillator. After recovering, Kerry Packer donated a large sum to the Ambulance Service of New South Wales in order that all ambulances in New South Wales should be fitted with a personal defibrillator, which is why defibrillators in Australia are sometimes colloquially called "Packer Whackers".
Gallery
| Biology and health sciences | Treatments | Health |
146393 | https://en.wikipedia.org/wiki/Mouth-to-mouth%20resuscitation | Mouth-to-mouth resuscitation | Mouth-to-mouth resuscitation, a form of artificial ventilation, is the act of assisting or stimulating respiration in which a rescuer presses their mouth against that of the victim and blows air into the person's lungs. Artificial respiration takes many forms, but generally entails providing air for a person who is not breathing or is not making sufficient respiratory effort on their own. It is used on a patient with a beating heart or as part of cardiopulmonary resuscitation (CPR) to achieve the internal respiration.
Pulmonary ventilation (and hence external respiration) is achieved through manual insufflation of the lungs either by the rescuer blowing into the patient's lungs, or by using a mechanical device to do so. This method of insufflation has been proved more effective than methods which involve mechanical manipulation of the patient's chest or arms, such as the Silvester method. It is also known as expired air resuscitation (EAR), expired air ventilation (EAV), rescue breathing, or colloquially the kiss of life. It was introduced as a life-saving measure in 1950.
Mouth-to-mouth resuscitation is a part of most protocols for performing cardiopulmonary resuscitation (CPR) making it an essential skill for first aid. In some situations, mouth-to-mouth resuscitation is also performed separately, for instance in near-drowning and opiate overdoses. The performance of mouth-to-mouth resuscitation on its own is now limited in most protocols to health professionals, whereas lay first-aiders are advised to undertake full CPR in any case where the patient is not breathing sufficiently.
History
In 1773, English physician William Hawes (1736–1808) began publicising the power of artificial respiration to resuscitate people who superficially appeared to have drowned. For a year he paid a reward out of his own pocket to any one bringing him a body rescued from the water within a reasonable time of immersion. Thomas Cogan, another English physician, who had become interested in the same subject during a stay at Amsterdam, where in 1767 a society for preservation of life from accidents in water was instituted, joined Hawes in his crusade. In the summer of 1774 Hawes and Cogan each brought fifteen friends to a meeting at the Chapter Coffee-house, St Paul's Churchyard, where they founded the Royal Humane Society as a campaigning group for first aid and resuscitation.
Gradually, branches of the Royal Humane Society were set up in other parts of the country, mainly in ports and coastal towns where the risk of drowning was high and by the end of the 19th century the society had upwards of 280 depots throughout the UK, supplied with life-saving apparatus. The earliest of these depots was the Receiving House in Hyde Park, on the north bank of the Serpentine, which was built in 1794 on a site granted by George III. Hyde Park was chosen because tens of thousands of people swam in the Serpentine in the summer and ice-skated in the winter. Boats and boatmen were kept to render aid to bathers, and in the winter ice-men were sent round to the different skating grounds in and around London.
The society distributed money-rewards, medals, clasps and testimonials, to those who saved or attempted to save drowning people. It further recognized "all cases of exceptional bravery in rescuing or attempting to rescue persons from asphyxia in mines, wells, blasting furnaces, or in sewers where foul gas may endanger life."
Insufflations
Insufflation, also known as 'rescue breaths' or 'ventilations', is the act of mechanically forcing air into a patient's respiratory system. This can be achieved via a number of methods, which will depend on the situation and equipment available. All methods require good airway management to perform, which ensures that the method is effective. These methods include:
Mouth-to-mouth - This involves the rescuer making a seal between his or her mouth and the patient's mouth and 'blowing', to pass air into the patient's body
Mouth-to-nose - In some instances, the rescuer may need or wish to form a seal with the patient's nose. Typical reasons for this include maxillofacial injuries, performing the procedure in water or the remains of vomit in the mouth
Mouth-to-face - Used on both animal muzzles and infants under 2, as this forms the most effective seal on both the mouth and nostrils
Mouth-to-mask – Most organisations recommend the use of some sort of barrier between rescuer and patient to reduce cross infection risk. One popular type is the 'pocket mask'. This may be able to provide higher tidal volumes than a Bag Valve Mask.
Adjuncts to insufflation
Most training organisations recommend that in any of the methods involving mouth-to-patient, that a protective barrier is used, to minimise the possibility of cross infection (in either direction).
Barriers available include pocket masks and keyring-sized face shields. These barriers are an example of personal protective equipment to guard the face against splashing, spraying or splattering of blood or other potentially infectious materials.
These barriers should provide a one-way filter valve which lets the air from the rescuer deliver to the patient while any substances from the patient (e.g. vomit, blood) cannot reach the rescuer. Many adjuncts are single use, though if they are multi use, after use of the adjunct, the mask must be cleaned and autoclaved and the filter replaced. It is very important for the mask to be replaced or cleaned because it can act as a transporter of various diseases.
The CPR mask is the preferred method of ventilating a patient when only one rescuer is available. Many feature inlets to support supplemental oxygen, which increases the oxygen being delivered from the approximate 17% available in the expired air of the rescuer to around 40-50%.
Efficiency of mouth-to-patient insufflation
Normal atmospheric air contains approximately 21% oxygen when inhaled. After gaseous exchange has taken place in the lungs, with waste products (notably carbon dioxide) moved from the bloodstream to the lungs, the air being exhaled by humans normally contains around 17% oxygen. This means that the human body utilises only around 19% of the oxygen inhaled, leaving over 80% of the oxygen available in the exhalatory breath.
This means that there is more than enough residual oxygen to be used in the lungs of the patient, which then enters the blood.
Oxygen
The efficiency of artificial respiration can be greatly increased by the simultaneous use of oxygen therapy. The amount of oxygen available to the patient in mouth-to-mouth is around 16%. If this is done through a pocket mask with an oxygen flow, this increases to 40% oxygen. If either a bag valve mask or a mechanical ventilator is used with an oxygen supply, this rises to 99% oxygen. The greater the oxygen concentration, the more efficient the gaseous exchange will be in the lungs.
| Biology and health sciences | Treatments | Health |
146396 | https://en.wikipedia.org/wiki/Tracheal%20intubation | Tracheal intubation | Tracheal intubation, usually simply referred to as intubation, is the placement of a flexible plastic tube into the trachea (windpipe) to maintain an open airway or to serve as a conduit through which to administer certain drugs. It is frequently performed in critically injured, ill, or anesthetized patients to facilitate ventilation of the lungs, including mechanical ventilation, and to prevent the possibility of asphyxiation or airway obstruction.
The most widely used route is orotracheal, in which an endotracheal tube is passed through the mouth and vocal apparatus into the trachea. In a nasotracheal procedure, an endotracheal tube is passed through the nose and vocal apparatus into the trachea. Other methods of intubation involve surgery and include the cricothyrotomy (used almost exclusively in emergency circumstances) and the tracheotomy, used primarily in situations where a prolonged need for airway support is anticipated.
Because it is an invasive and uncomfortable medical procedure, intubation is usually performed after administration of general anesthesia and a neuromuscular-blocking drug. It can, however, be performed in the awake patient with local or topical anesthesia or in an emergency without any anesthesia at all. Intubation is normally facilitated by using a conventional laryngoscope, flexible fiberoptic bronchoscope, or video laryngoscope to identify the vocal cords and pass the tube between them into the trachea instead of into the esophagus. Other devices and techniques may be used alternatively.
After the trachea has been intubated, a balloon cuff is typically inflated just above the far end of the tube to help secure it in place, to prevent leakage of respiratory gases, and to protect the tracheobronchial tree from receiving undesirable material such as stomach acid. The tube is then secured to the face or neck and connected to a T-piece, anesthesia breathing circuit, bag valve mask device, or a mechanical ventilator. Once there is no longer a need for ventilatory assistance or protection of the airway, the tracheal tube is removed; this is referred to as extubation of the trachea (or decannulation, in the case of a surgical airway such as a cricothyrotomy or a tracheotomy).
For centuries, tracheotomy was considered the only reliable method for intubation of the trachea. However, because only a minority of patients survived the operation, physicians undertook tracheotomy only as a last resort, on patients who were nearly dead. It was not until the late 19th century, however, that advances in understanding of anatomy and physiology, as well an appreciation of the germ theory of disease, had improved the outcome of this operation to the point that it could be considered an acceptable treatment option. Also at that time, advances in endoscopic instrumentation had improved to such a degree that direct laryngoscopy had become a viable means to secure the airway by the non-surgical orotracheal route. By the mid-20th century, the tracheotomy as well as endoscopy and non-surgical tracheal intubation had evolved from rarely employed procedures to becoming essential components of the practices of anesthesiology, critical care medicine, emergency medicine, and laryngology.
Tracheal intubation can be associated with complications such as broken teeth or lacerations of the tissues of the upper airway. It can also be associated with potentially fatal complications such as pulmonary aspiration of stomach contents which can result in a severe and sometimes fatal chemical aspiration pneumonitis, or unrecognized intubation of the esophagus which can lead to potentially fatal anoxia. Because of this, the potential for difficulty or complications due to the presence of unusual airway anatomy or other uncontrolled variables is carefully evaluated before undertaking tracheal intubation. Alternative strategies for securing the airway must always be readily available.
Indications
Tracheal intubation is indicated in a variety of situations when illness or a medical procedure prevents a person from maintaining a clear airway, breathing, and oxygenating the blood. In these circumstances, oxygen supplementation using a simple face mask is inadequate.
Depressed level of consciousness
Perhaps the most common indication for tracheal intubation is for the placement of a conduit through which nitrous oxide or volatile anesthetics may be administered. General anesthetic agents, opioids, and neuromuscular-blocking drugs may diminish or even abolish the respiratory drive. Although it is not the only means to maintain a patent airway during general anesthesia, intubation of the trachea provides the most reliable means of oxygenation and ventilation and the greatest degree of protection against regurgitation and pulmonary aspiration.
Damage to the brain (such as from a massive stroke, non-penetrating head injury, intoxication or poisoning) may result in a depressed level of consciousness. When this becomes severe to the point of stupor or coma (defined as a score on the Glasgow Coma Scale of less than 8), dynamic collapse of the extrinsic muscles of the airway can obstruct the airway, impeding the free flow of air into the lungs. Furthermore, protective airway reflexes such as coughing and swallowing may be diminished or absent. Tracheal intubation is often required to restore patency (the relative absence of blockage) of the airway and protect the tracheobronchial tree from pulmonary aspiration of gastric contents.
Hypoxemia
Intubation may be necessary for a patient with decreased oxygen content and oxygen saturation of the blood caused when their breathing is inadequate (hypoventilation), suspended (apnea), or when the lungs are unable to sufficiently transfer gasses to the blood. Such patients, who may be awake and alert, are typically critically ill with a multisystem disease or multiple severe injuries. Examples of such conditions include cervical spine injury, multiple rib fractures, severe pneumonia, acute respiratory distress syndrome (ARDS), or near-drowning. Specifically, intubation is considered if the arterial partial pressure of oxygen (PaO2) is less than 60 millimeters of mercury (mm Hg) while breathing an inspired O2 concentration (FIO2) of 50% or greater. In patients with elevated arterial carbon dioxide, an arterial partial pressure of CO2 (PaCO2) greater than 45 mm Hg in the setting of acidemia would prompt intubation, especially if a series of measurements demonstrate a worsening respiratory acidosis. Regardless of the laboratory values, these guidelines are always interpreted in the clinical context.
Airway obstruction
Actual or impending airway obstruction is a common indication for intubation of the trachea. Life-threatening airway obstruction may occur when a foreign body becomes lodged in the airway; this is especially common in infants and toddlers. Severe blunt or penetrating injury to the face or neck may be accompanied by swelling and an expanding hematoma, or injury to the larynx, trachea or bronchi. Airway obstruction is also common in people who have suffered smoke inhalation or burns within or near the airway or epiglottitis. Sustained generalized seizure activity and angioedema are other common causes of life-threatening airway obstruction which may require tracheal intubation to secure the airway.
Manipulation of the airway
Diagnostic or therapeutic manipulation of the airway (such as bronchoscopy, laser therapy or stenting of the bronchi) may intermittently interfere with the ability to breathe; intubation may be necessary in such situations.
Newborns
Syndromes such as respiratory distress syndrome, congenital heart disease, pneumothorax, and shock may lead to breathing problems in newborn infants that require endotracheal intubation and mechanically assisted breathing (mechanical ventilation). Newborn infants may also require endotracheal intubation during surgery while under general anaesthesia.
Equipment
Laryngoscopes
The vast majority of tracheal intubations involve the use of a viewing instrument of one type or another. The modern conventional laryngoscope consists of a handle containing batteries that power a light and a set of interchangeable blades, which are either straight or curved. This device is designed to allow the laryngoscopist to directly view the larynx. Due to the widespread availability of such devices, the technique of blind intubation of the trachea is rarely practiced today, although it may still be useful in certain emergency situations, such as natural or man-made disasters. In the prehospital emergency setting, digital intubation may be necessitated if the patient is in a position that makes direct laryngoscopy impossible. For example, digital intubation may be used by a paramedic if the patient is entrapped in an inverted position in a vehicle after a motor vehicle collision with a prolonged extrication time.
The decision to use a straight or curved laryngoscope blade depends partly on the specific anatomical features of the airway, and partly on the personal experience and preference of the laryngoscopist. The Miller blade, characterized by its straight, elongated shape with a curved tip, is frequently employed in patients with challenging airway anatomy, such as those with limited mouth opening or a high larynx. Its design allows for direct visualization of the epiglottis, facilitating precise glottic exposure.
Conversely, the Macintosh blade, with its curved configuration reminiscent of the letters "C" or "J," is favored in routine intubations for patients with normal airway anatomy. Its curved design enables indirect laryngoscopy, providing enhanced visualization of the vocal cords and glottis in most adult patients.
The choice between the Miller and Macintosh blades is influenced by specific anatomical considerations and the preferences of the laryngoscopist. While the Macintosh blade is the most commonly utilized curved laryngoscope blade, the Miller blade is the preferred option for straight blade intubation. Both blades are available in various sizes, ranging from size 0 (infant) to size 4 (large adult), catering to patients of different ages and anatomies. Additionally, there exists a myriad of specialty blades with unique features, including mirrors for enhanced visualization and ports for oxygen administration, primarily utilized by anesthetists and otolaryngologists in operating room settings.
Fiberoptic laryngoscopes have become increasingly available since the 1990s. In contrast to the conventional laryngoscope, these devices allow the laryngoscopist to indirectly view the larynx. This provides a significant advantage in situations where the operator needs to see around an acute bend in order to visualize the glottis, and deal with otherwise difficult intubations. Video laryngoscopes are specialized fiberoptic laryngoscopes that use a digital video camera sensor to allow the operator to view the glottis and larynx on a video monitor. Other "noninvasive" devices which can be employed to assist in tracheal intubation are the laryngeal mask airway (used as a conduit for endotracheal tube placement) and the Airtraq.
Stylets
An intubating stylet is a malleable metal wire designed to be inserted into the endotracheal tube to make the tube conform better to the upper airway anatomy of the specific individual. This aid is commonly used with a difficult laryngoscopy. Just as with laryngoscope blades, there are also several types of available stylets, such as the Verathon Stylet, which is specifically designed to follow the 60° blade angle of the GlideScope video laryngoscope.
The Eschmann tracheal tube introducer (also referred to as a "gum elastic bougie") is specialized type of stylet used to facilitate difficult intubation. This flexible device is in length, 15 French (5 mm diameter) with a small "hockey-stick" angle at the far end. Unlike a traditional intubating stylet, the Eschmann tracheal tube introducer is typically inserted directly into the trachea and then used as a guide over which the endotracheal tube can be passed (in a manner analogous to the Seldinger technique). As the Eschmann tracheal tube introducer is considerably less rigid than a conventional stylet, this technique is considered to be a relatively atraumatic means of tracheal intubation.
The tracheal tube exchanger is a hollow catheter, in length, that can be used for removal and replacement of tracheal tubes without the need for laryngoscopy. The Cook Airway Exchange Catheter (CAEC) is another example of this type of catheter; this device has a central lumen (hollow channel) through which oxygen can be administered. Airway exchange catheters are long hollow catheters which often have connectors for jet ventilation, manual ventilation, or oxygen insufflation. It is also possible to connect the catheter to a capnograph to perform respiratory monitoring.
The lighted stylet is a device that employs the principle of transillumination to facilitate blind orotracheal intubation (an intubation technique in which the laryngoscopist does not view the glottis).
Tracheal tubes
A tracheal tube is a catheter that is inserted into the trachea for the primary purpose of establishing and maintaining a patent (open and unobstructed) airway. Tracheal tubes are frequently used for airway management in the settings of general anesthesia, critical care, mechanical ventilation, and emergency medicine. Many different types of tracheal tubes are available, suited for different specific applications. An endotracheal tube is a specific type of tracheal tube that is nearly always inserted through the mouth (orotracheal) or nose (nasotracheal). It is a breathing conduit designed to be placed into the airway of critically injured, ill or anesthetized patients in order to perform mechanical positive pressure ventilation of the lungs and to prevent the possibility of aspiration or airway obstruction. The endotracheal tube has a fitting designed to be connected to a source of pressurized gas such as oxygen. At the other end is an orifice through which such gases are directed into the lungs and may also include a balloon (referred to as a cuff). The tip of the endotracheal tube is positioned above the carina (before the trachea divides to each lung) and sealed within the trachea so that the lungs can be ventilated equally. A tracheostomy tube is another type of tracheal tube; this curved metal or plastic tube is inserted into a tracheostomy stoma or a cricothyrotomy incision.
Tracheal tubes can be used to ensure the adequate exchange of oxygen and carbon dioxide, to deliver oxygen in higher concentrations than found in air, or to administer other gases such as helium, nitric oxide, nitrous oxide, xenon, or certain volatile anesthetic agents such as desflurane, isoflurane, or sevoflurane. They may also be used as a route for administration of certain medications such as bronchodilators, inhaled corticosteroids, and drugs used in treating cardiac arrest such as atropine, epinephrine, lidocaine and vasopressin.
Originally made from latex rubber, most modern endotracheal tubes today are constructed of polyvinyl chloride. Tubes constructed of silicone rubber, wire-reinforced silicone rubber or stainless steel are also available for special applications. For human use, tubes range in size from in internal diameter. The size is chosen based on the patient's body size, with the smaller sizes being used for infants and children. Most endotracheal tubes have an inflatable cuff to seal the tracheobronchial tree against leakage of respiratory gases and pulmonary aspiration of gastric contents, blood, secretions, and other fluids. Uncuffed tubes are also available, though their use is limited mostly to children (in small children, the cricoid cartilage is the narrowest portion of the airway and usually provides an adequate seal for mechanical ventilation).
In addition to cuffed or uncuffed, preformed endotracheal tubes are also available. The oral and nasal RAE tubes (named after the inventors Ring, Adair and Elwyn) are the most widely used of the preformed tubes.
There are a number of different types of double-lumen endo-bronchial tubes that have endobronchial as well as endotracheal channels (Carlens, White and Robertshaw tubes). These tubes are typically coaxial, with two separate channels and two separate openings. They incorporate an endotracheal lumen which terminates in the trachea and an endobronchial lumen, the distal tip of which is positioned 1–2 cm into the right or left mainstem bronchus. There is also the Univent tube, which has a single tracheal lumen and an integrated endobronchial blocker. These tubes enable one to ventilate both lungs, or either lung independently. Single-lung ventilation (allowing the lung on the operative side to collapse) can be useful during thoracic surgery, as it can facilitate the surgeon's view and access to other relevant structures within the thoracic cavity.
The "armored" endotracheal tubes are cuffed, wire-reinforced silicone rubber tubes. They are much more flexible than polyvinyl chloride tubes, yet they are difficult to compress or kink. This can make them useful for situations in which the trachea is anticipated to remain intubated for a prolonged duration, or if the neck is to remain flexed during surgery. Most armored tubes have a Magill curve, but preformed armored RAE tubes are also available. Another type of endotracheal tube has four small openings just above the inflatable cuff, which can be used for suction of the trachea or administration of intratracheal medications if necessary. Other tubes (such as the Bivona Fome-Cuf tube) are designed specifically for use in laser surgery in and around the airway.
Methods to confirm tube placement
No single method for confirming tracheal tube placement has been shown to be 100% reliable. Accordingly, the use of multiple methods for confirmation of correct tube placement is now widely considered to be the standard of care. Such methods include direct visualization as the tip of the tube passes through the glottis, or indirect visualization of the tracheal tube within the trachea using a device such as a bronchoscope. With a properly positioned tracheal tube, equal bilateral breath sounds will be heard upon listening to the chest with a stethoscope, and no sound upon listening to the area over the stomach. Equal bilateral rise and fall of the chest wall will be evident with ventilatory excursions. A small amount of water vapor will also be evident within the lumen of the tube with each exhalation and there will be no gastric contents in the tracheal tube at any time.
Ideally, at least one of the methods utilized for confirming tracheal tube placement will be a measuring instrument. Waveform capnography has emerged as the gold standard for the confirmation of tube placement within the trachea. Other methods relying on instruments include the use of a colorimetric end-tidal carbon dioxide detector, a self-inflating esophageal bulb, or an esophageal detection device. The distal tip of a properly positioned tracheal tube will be located in the mid-trachea, roughly above the bifurcation of the carina; this can be confirmed by chest x-ray. If it is inserted too far into the trachea (beyond the carina), the tip of the tracheal tube is likely to be within the right main bronchus—a situation often referred to as a "right mainstem intubation". In this situation, the left lung may be unable to participate in ventilation, which can lead to decreased oxygen content due to ventilation/perfusion mismatch.
Special situations
Emergencies
Tracheal intubation in the emergency setting can be difficult with the fiberoptic bronchoscope due to blood, vomit, or secretions in the airway and poor patient cooperation. Because of this, patients with massive facial injury, complete upper airway obstruction, severely diminished ventilation, or profuse upper airway bleeding are poor candidates for fiberoptic intubation. Fiberoptic intubation under general anesthesia typically requires two skilled individuals. Success rates of only 83–87% have been reported using fiberoptic techniques in the emergency department, with significant nasal bleeding occurring in up to 22% of patients. These drawbacks limit the use of fiberoptic bronchoscopy somewhat in urgent and emergency situations.
Personnel experienced in direct laryngoscopy are not always immediately available in certain settings that require emergency tracheal intubation. For this reason, specialized devices have been designed to act as bridges to a definitive airway. Such devices include the laryngeal mask airway, cuffed oropharyngeal airway and the esophageal-tracheal combitube (Combitube). Other devices such as rigid stylets, the lightwand (a blind technique) and indirect fiberoptic rigid stylets, such as the Bullard scope, Upsher scope and the WuScope can also be used as alternatives to direct laryngoscopy. Each of these devices have its own unique set of benefits and drawbacks, and none of them is effective under all circumstances.
Rapid-sequence induction and intubation
Rapid sequence induction and intubation (RSI) is a particular method of induction of general anesthesia, commonly employed in emergency operations and other situations where patients are assumed to have a full stomach. The objective of RSI is to minimize the possibility of regurgitation and pulmonary aspiration of gastric contents during the induction of general anesthesia and subsequent tracheal intubation. RSI traditionally involves preoxygenating the lungs with a tightly fitting oxygen mask, followed by the sequential administration of an intravenous sleep-inducing agent and a rapidly acting neuromuscular-blocking drug, such as rocuronium, succinylcholine, or cisatracurium besilate, before intubation of the trachea.
One important difference between RSI and routine tracheal intubation is that the practitioner does not manually assist the ventilation of the lungs after the onset of general anesthesia and cessation of breathing, until the trachea has been intubated and the cuff has been inflated. Another key feature of RSI is the application of manual 'cricoid pressure' to the cricoid cartilage, often referred to as the "Sellick maneuver", prior to instrumentation of the airway and intubation of the trachea.
Named for British anesthetist Brian Arthur Sellick (1918–1996) who first described the procedure in 1961, the goal of cricoid pressure is to minimize the possibility of regurgitation and pulmonary aspiration of gastric contents. Cricoid pressure has been widely used during RSI for nearly fifty years, despite a lack of compelling evidence to support this practice. The initial article by Sellick was based on a small sample size at a time when high tidal volumes, head-down positioning and barbiturate anesthesia were the rule. Beginning around 2000, a significant body of evidence has accumulated which questions the effectiveness of cricoid pressure. The application of cricoid pressure may in fact displace the esophagus laterally instead of compressing it as described by Sellick. Cricoid pressure may also compress the glottis, which can obstruct the view of the laryngoscopist and actually cause a delay in securing the airway.
Cricoid pressure is often confused with the "BURP" (Backwards Upwards Rightwards Pressure) maneuver. While both of these involve digital pressure to the anterior aspect (front) of the laryngeal apparatus, the purpose of the latter is to improve the view of the glottis during laryngoscopy and tracheal intubation, rather than to prevent regurgitation. Both cricoid pressure and the BURP maneuver have the potential to worsen laryngoscopy.
RSI may also be used in prehospital emergency situations when a patient is conscious but respiratory failure is imminent (such as in extreme trauma). This procedure is commonly performed by flight paramedics. Flight paramedics often use RSI to intubate before transport because intubation in a moving fixed-wing or rotary-wing aircraft is extremely difficult to perform due to environmental factors. The patient will be paralyzed and intubated on the ground before transport by aircraft.
Cricothyrotomy
A cricothyrotomy is an incision made through the skin and cricothyroid membrane to establish a patent airway during certain life-threatening situations, such as airway obstruction by a foreign body, angioedema, or massive facial trauma. A cricothyrotomy is nearly always performed as a last resort in cases where orotracheal and nasotracheal intubation are impossible or contraindicated. Cricothyrotomy is easier and quicker to perform than tracheotomy, does not require manipulation of the cervical spine and is associated with fewer complications.
The easiest method to perform this technique is the needle cricothyrotomy (also referred to as a percutaneous dilational cricothyrotomy), in which a large-bore (12–14 gauge) intravenous catheter is used to puncture the cricothyroid membrane. Oxygen can then be administered through this catheter via jet insufflation. However, while needle cricothyrotomy may be life-saving in extreme circumstances, this technique is only intended to be a temporizing measure until a definitive airway can be established. While needle cricothyrotomy can provide adequate oxygenation, the small diameter of the cricothyrotomy catheter is insufficient for elimination of carbon dioxide (ventilation). After one hour of apneic oxygenation through a needle cricothyrotomy, one can expect a PaCO2 of greater than 250 mm Hg and an arterial pH of less than 6.72, despite an oxygen saturation of 98% or greater. A more definitive airway can be established by performing a surgical cricothyrotomy, in which a endotracheal tube or tracheostomy tube can be inserted through a larger incision.
Several manufacturers market prepackaged cricothyrotomy kits, which enable one to use either a wire-guided percutaneous dilational (Seldinger) technique, or the classic surgical technique to insert a polyvinylchloride catheter through the cricothyroid membrane. The kits may be stocked in hospital emergency departments and operating suites, as well as ambulances and other selected pre-hospital settings.
Tracheotomy
Tracheotomy consists of making an incision on the front of the neck and opening a direct airway through an incision in the trachea. The resulting opening can serve independently as an airway or as a site for a tracheostomy tube to be inserted; this tube allows a person to breathe without the use of his nose or mouth. The opening may be made by a scalpel or a needle (referred to as surgical and percutaneous techniques respectively) and both techniques are widely used in current practice. In order to limit the risk of damage to the recurrent laryngeal nerves (the nerves that control the voice box), the tracheotomy is performed as high in the trachea as possible. If only one of these nerves is damaged, the patient's voice may be impaired (dysphonia); if both of the nerves are damaged, the patient will be unable to speak (aphonia). In the acute setting, indications for tracheotomy are similar to those for cricothyrotomy. In the chronic setting, indications for tracheotomy include the need for long-term mechanical ventilation and removal of tracheal secretions (e.g., comatose patients, or extensive surgery involving the head and neck).
Children
There are significant differences in airway anatomy and respiratory physiology between children and adults, and these are taken into careful consideration before performing tracheal intubation of any pediatric patient. The differences, which are quite significant in infants, gradually disappear as the human body approaches a mature age and body mass index.
For infants and young children, orotracheal intubation is easier than the nasotracheal route. Nasotracheal intubation carries a risk of dislodgement of adenoids and nasal bleeding. Despite the greater difficulty, nasotracheal intubation route is preferable to orotracheal intubation in children undergoing intensive care and requiring prolonged intubation because this route allows a more secure fixation of the tube. As with adults, there are a number of devices specially designed for assistance with difficult tracheal intubation in children. Confirmation of proper position of the tracheal tube is accomplished as with adult patients.
Because the airway of a child is narrow, a small amount of glottic or tracheal swelling can produce critical obstruction. Inserting a tube that is too large relative to the diameter of the trachea can cause swelling. Conversely, inserting a tube that is too small can result in inability to achieve effective positive pressure ventilation due to retrograde escape of gas through the glottis and out the mouth and nose (often referred to as a "leak" around the tube). An excessive leak can usually be corrected by inserting a larger tube or a cuffed tube.
The tip of a correctly positioned tracheal tube will be in the mid-trachea, between the collarbones on an anteroposterior chest radiograph. The correct diameter of the tube is that which results in a small leak at a pressure of about of water. The appropriate inner diameter for the endotracheal tube is estimated to be roughly the same diameter as the child's little finger. The appropriate length for the endotracheal tube can be estimated by doubling the distance from the corner of the child's mouth to the ear canal. For premature infants internal diameter is an appropriate size for the tracheal tube. For infants of normal gestational age, internal diameter is an appropriate size. For normally nourished children 1 year of age and older, two formulae are used to estimate the appropriate diameter and depth for tracheal intubation. The internal diameter of the tube in mm is (patient's age in years + 16) / 4, while the appropriate depth of insertion in cm is 12 + (patient's age in years / 2).
Newborn infants
Endotrachael suctioning is often used during intubation in newborn infants to reduce the risk of a blocked tube due to secretions, a collapsed lung, and to reduce pain. Suctioning is sometimes used at specifically scheduled intervals, "as needed", and less frequently. Further research is necessary to determine the most effective suctioning schedule or frequency of suctioning in intubated infants.
In newborns free flow oxygen used to be recommended during intubation however as there is no evidence of benefit the 2011 NRP guidelines no longer do.
Predicting difficulty
Tracheal intubation is not a simple procedure and the consequences of failure are grave. Therefore, the patient is carefully evaluated for potential difficulty or complications beforehand. This involves taking the medical history of the patient and performing a physical examination, the results of which can be scored against one of several classification systems. The proposed surgical procedure (e.g., surgery involving the head and neck, or bariatric surgery) may lead one to anticipate difficulties with intubation. Many individuals have unusual airway anatomy, such as those who have limited movement of their neck or jaw, or those who have tumors, deep swelling due to injury or to allergy, developmental abnormalities of the jaw, or excess fatty tissue of the face and neck. Using conventional laryngoscopic techniques, intubation of the trachea can be difficult or even impossible in such patients. This is why all persons performing tracheal intubation must be familiar with alternative techniques of securing the airway. Use of the flexible fiberoptic bronchoscope and similar devices has become among the preferred techniques in the management of such cases. However, these devices require a different skill set than that employed for conventional laryngoscopy and are expensive to purchase, maintain and repair.
When taking the patient's medical history, the subject is questioned about any significant signs or symptoms, such as difficulty in speaking or difficulty in breathing. These may suggest obstructing lesions in various locations within the upper airway, larynx, or tracheobronchial tree. A history of previous surgery (e.g., previous cervical fusion), injury, radiation therapy, or tumors involving the head, neck and upper chest can also provide clues to a potentially difficult intubation. Previous experiences with tracheal intubation, especially difficult intubation, intubation for prolonged duration (e.g., intensive care unit) or prior tracheotomy are also noted.
A detailed physical examination of the airway is important, particularly:
the range of motion of the cervical spine: the subject should be able to tilt the head back and then forward so that the chin touches the chest.
the range of motion of the jaw (the temporomandibular joint): three of the subject's fingers should be able to fit between the upper and lower incisors.
the size and shape of the upper jaw and lower jaw, looking especially for problems such as maxillary hypoplasia (an underdeveloped upper jaw), micrognathia (an abnormally small jaw), or retrognathia (misalignment of the upper and lower jaw).
the thyromental distance: three of the subject's fingers should be able to fit between the Adam's apple and the chin.
the size and shape of the tongue and palate relative to the size of the mouth.
the teeth, especially noting the presence of prominent maxillary incisors, any loose or damaged teeth, or crowns.
Many classification systems have been developed in an effort to predict difficulty of tracheal intubation, including the Cormack-Lehane classification system, the Intubation Difficulty Scale (IDS), and the Mallampati score. The Mallampati score is drawn from the observation that the size of the base of the tongue influences the difficulty of intubation. It is determined by looking at the anatomy of the mouth, and in particular the visibility of the base of palatine uvula, faucial pillars and the soft palate. Although such medical scoring systems may aid in the evaluation of patients, no single score or combination of scores can be trusted to specifically detect all and only those patients who are difficult to intubate. Furthermore, one study of experienced anesthesiologists, on the widely used Cormack–Lehane classification system, found they did not score the same patients consistently over time, and that only 25% could correctly define all four grades of the widely used Cormack–Lehane classification system. Under certain emergency circumstances (e.g., severe head trauma or suspected cervical spine injury), it may be impossible to fully utilize these the physical examination and the various classification systems to predict the difficulty of tracheal intubation. A Cochrane systematic review examined the sensitivity and specificity of various bedside tests commonly used for predicting difficulty in airway management. In such cases, alternative techniques of securing the airway must be readily available.
Complications
Tracheal intubation is generally considered the best method for airway management under a wide variety of circumstances, as it provides the most reliable means of oxygenation and ventilation and the greatest degree of protection against regurgitation and pulmonary aspiration. However, tracheal intubation requires a great deal of clinical experience to master and serious complications may result even when properly performed.
Four anatomic features must be present for orotracheal intubation to be straightforward: adequate mouth opening (full range of motion of the temporomandibular joint), sufficient pharyngeal space (determined by examining the back of the mouth), sufficient submandibular space (distance between the thyroid cartilage and the chin, the space into which the tongue must be displaced in order for the larygoscopist to view the glottis), and adequate extension of the cervical spine at the atlanto-occipital joint. If any of these variables is in any way compromised, intubation should be expected to be difficult.
Minor complications are common after laryngoscopy and insertion of an orotracheal tube. These are typically of short duration, such as sore throat, lacerations of the lips or gums or other structures within the upper airway, chipped, fractured or dislodged teeth, and nasal injury. Other complications which are common but potentially more serious include accelerated or irregular heartbeat, high blood pressure, elevated intracranial and introcular pressure, and bronchospasm.
More serious complications include laryngospasm, perforation of the trachea or esophagus, pulmonary aspiration of gastric contents or other foreign bodies, fracture or dislocation of the cervical spine, temporomandibular joint or arytenoid cartilages, decreased oxygen content, elevated arterial carbon dioxide, and vocal cord weakness. In addition to these complications, tracheal intubation via the nasal route carries a risk of dislodgement of adenoids and potentially severe nasal bleeding. Newer technologies such as flexible fiberoptic laryngoscopy have fared better in reducing the incidence of some of these complications, though the most frequent cause of intubation trauma remains a lack of skill on the part of the laryngoscopist.
Complications may also be severe and long-lasting or permanent, such as vocal cord damage, esophageal perforation and retropharyngeal abscess, bronchial intubation, or nerve injury. They may even be immediately life-threatening, such as laryngospasm and negative pressure pulmonary edema (fluid in the lungs), aspiration, unrecognized esophageal intubation, or accidental disconnection or dislodgement of the tracheal tube. Potentially fatal complications more often associated with prolonged intubation or tracheotomy include abnormal communication between the trachea and nearby structures such as the innominate artery (tracheoinnominate fistula) or esophagus (tracheoesophageal fistula). Other significant complications include airway obstruction due to loss of tracheal rigidity, ventilator-associated pneumonia and narrowing of the glottis or trachea. The cuff pressure is monitored carefully in order to avoid complications from over-inflation, many of which can be traced to excessive cuff pressure restricting the blood supply to the tracheal mucosa. A 2000 Spanish study of bedside percutaneous tracheotomy reported overall complication rates of 10–15% and procedural mortality of 0%, which is comparable to those of other series reported in the literature from the Netherlands and the United States.
Inability to secure the airway, with subsequent failure of oxygenation and ventilation is a life-threatening complication which if not immediately corrected leads to decreased oxygen content, brain damage, cardiovascular collapse, and death. When performed improperly, the associated complications (e.g., unrecognized esophageal intubation) may be rapidly fatal. Without adequate training and experience, the incidence of such complications is high. The case of Andrew Davis Hughes, from Emerald Isle, NC is a widely known case in which the patient was improperly intubated and, due to the lack of oxygen, sustained severe brain damage and died. For example, among paramedics in several United States urban communities, unrecognized esophageal or hypopharyngeal intubation has been reported to be 6% to 25%. Although not common, where basic emergency medical technicians are permitted to intubate, reported success rates are as low as 51%. In one study, nearly half of patients with misplaced tracheal tubes died in the emergency room. Because of this, the American Heart Association's Guidelines for Cardiopulmonary Resuscitation have de-emphasized the role of tracheal intubation in favor of other airway management techniques such as bag-valve-mask ventilation, the laryngeal mask airway and the Combitube. Higher quality studies demonstrate favorable evidence for this shift, as they have shown no survival or neurological benefit with endotracheal intubation over supraglottic airway devices (Laryngeal mask or Combitube).
One complication—unintentional and unrecognized intubation of the esophagus—is both common (as frequent as 25% in the hands of inexperienced personnel) and likely to result in a deleterious or even fatal outcome. In such cases, oxygen is inadvertently administered to the stomach, from where it cannot be taken up by the circulatory system, instead of the lungs. If this situation is not immediately identified and corrected, death will ensue from cerebral and cardiac anoxia.
Of 4,460 claims in the American Society of Anesthesiologists (ASA) Closed Claims Project database, 266 (approximately 6%) were for airway injury. Of these 266 cases, 87% of the injuries were temporary, 5% were permanent or disabling, and 8% resulted in death. Difficult intubation, age older than 60 years, and female gender were associated with claims for perforation of the esophagus or pharynx. Early signs of perforation were present in only 51% of perforation claims, whereas late sequelae occurred in 65%.
During the SARS and COVID-19 pandemics, tracheal intubation has been used with a ventilator in severe cases where the patient struggles to breathe. Performing the procedure carries a risk of the caregiver becoming infected.
Alternatives
Although it offers the greatest degree of protection against regurgitation and pulmonary aspiration, tracheal intubation is not the only means to maintain a patent airway. Alternative techniques for airway management and delivery of oxygen, volatile anesthetics or other breathing gases include the laryngeal mask airway, i-gel, cuffed oropharyngeal airway, continuous positive airway pressure (CPAP mask), nasal BiPAP mask, simple face mask, and nasal cannula.
General anesthesia is often administered without tracheal intubation in selected cases where the procedure is brief in duration, or procedures where the depth of anesthesia is not sufficient to cause significant compromise in ventilatory function. Even for longer duration or more invasive procedures, a general anesthetic may be administered without intubating the trachea, provided that patients are carefully selected, and the risk-benefit ratio is favorable (i.e., the risks associated with an unprotected airway are believed to be less than the risks of intubating the trachea).
Airway management can be classified into closed or open techniques depending on the system of ventilation used. Tracheal intubation is a typical example of a closed technique as ventilation occurs using a closed circuit. Several open techniques exist, such as spontaneous ventilation, apnoeic ventilation or jet ventilation. Each has its own specific advantages and disadvantages which determine when it should be used.
Spontaneous ventilation has been traditionally performed with an inhalational agent (i.e. gas induction or inhalational induction using halothane or sevoflurane) however it can also be performed using intravenous anaesthesia (e.g. propofol, ketamine or dexmedetomidine). SponTaneous Respiration using IntraVEnous anaesthesia and High-flow nasal oxygen (STRIVE Hi) is an open airway technique that uses an upwards titration of propofol which maintains ventilation at deep levels of anaesthesia. It has been used in airway surgery as an alternative to tracheal intubation.
History
Tracheotomy
The earliest known depiction of a tracheotomy is found on two Egyptian tablets dating back to around 3600 BC. The 110-page Ebers Papyrus, an Egyptian medical papyrus which dates to roughly 1550 BC, also makes reference to the tracheotomy. Tracheotomy was described in the Rigveda, a Sanskrit text of ayurvedic medicine written around 2000 BC in ancient India. The Sushruta Samhita from around 400 BC is another text from the Indian subcontinent on ayurvedic medicine and surgery that mentions tracheotomy. Asclepiades of Bithynia (–40 BC) is often credited as being the first physician to perform a non-emergency tracheotomy. Galen of Pergamon (AD 129–199) clarified the anatomy of the trachea and was the first to demonstrate that the larynx generates the voice. In one of his experiments, Galen used bellows to inflate the lungs of a dead animal. Ibn Sīnā (980–1037) described the use of tracheal intubation to facilitate breathing in 1025 in his 14-volume medical encyclopedia, The Canon of Medicine. In the 12th century medical textbook Al-Taisir, Ibn Zuhr (1092–1162)—also known as Avenzoar—of Al-Andalus provided a correct description of the tracheotomy operation.
The first detailed descriptions of tracheal intubation and subsequent artificial respiration of animals were from Andreas Vesalius (1514–1564) of Brussels. In his landmark book published in 1543, De humani corporis fabrica, he described an experiment in which he passed a reed into the trachea of a dying animal whose thorax had been opened and maintained ventilation by blowing into the reed intermittently. Antonio Musa Brassavola (1490–1554) of Ferrara successfully treated a patient with peritonsillar abscess by tracheotomy. Brassavola published his account in 1546; this operation has been identified as the first recorded successful tracheotomy, despite the many previous references to this operation. Towards the end of the 16th century, Hieronymus Fabricius (1533–1619) described a useful technique for tracheotomy in his writings, although he had never actually performed the operation himself. In 1620 the French surgeon Nicholas Habicot (1550–1624) published a report of four successful tracheotomies. In 1714, anatomist Georg Detharding (1671–1747) of the University of Rostock performed a tracheotomy on a drowning victim.
Despite the many recorded instances of its use since antiquity, it was not until the early 19th century that the tracheotomy finally began to be recognized as a legitimate means of treating severe airway obstruction. In 1852, French physician Armand Trousseau (1801–1867) presented a series of 169 tracheotomies to the Académie Impériale de Médecine. 158 of these were performed for the treatment of croup, and 11 were performed for "chronic maladies of the larynx". Between 1830 and 1855, more than 350 tracheotomies were performed in Paris, most of them at the Hôpital des Enfants Malades, a public hospital, with an overall survival rate of only 20–25%. This compares with 58% of the 24 patients in Trousseau's private practice, who fared better due to greater postoperative care.
In 1871, the German surgeon Friedrich Trendelenburg (1844–1924) published a paper describing the first successful elective human tracheotomy to be performed for the purpose of administration of general anesthesia. In 1888, Sir Morell Mackenzie (1837–1892) published a book discussing the indications for tracheotomy. In the early 20th century, tracheotomy became a life-saving treatment for patients affected with paralytic poliomyelitis who required mechanical ventilation. In 1909, Philadelphia laryngologist Chevalier Jackson (1865–1958) described a technique for tracheotomy that is used to this day.
Laryngoscopy and non-surgical techniques
In 1854, a Spanish singing teacher named Manuel García (1805–1906) became the first man to view the functioning glottis in a living human. In 1858, French pediatrician Eugène Bouchut (1818–1891) developed a new technique for non-surgical orotracheal intubation to bypass laryngeal obstruction resulting from a diphtheria-related pseudomembrane. In 1880, Scottish surgeon William Macewen (1848–1924) reported on his use of orotracheal intubation as an alternative to tracheotomy to allow a patient with glottic edema to breathe, as well as in the setting of general anesthesia with chloroform. In 1895, Alfred Kirstein (1863–1922) of Berlin first described direct visualization of the vocal cords, using an esophagoscope he had modified for this purpose; he called this device an autoscope.
In 1913, Chevalier Jackson was the first to report a high rate of success for the use of direct laryngoscopy as a means to intubate the trachea. Jackson introduced a new laryngoscope blade that incorporated a component that the operator could slide out to allow room for passage of an endotracheal tube or bronchoscope. Also in 1913, New York surgeon Henry H. Janeway (1873–1921) published results he had achieved using a laryngoscope he had recently developed. Another pioneer in this field was Sir Ivan Whiteside Magill (1888–1986), who developed the technique of awake blind nasotracheal intubation, the Magill forceps, the Magill laryngoscope blade, and several apparati for the administration of volatile anesthetic agents. The Magill curve of an endotracheal tube is also named for Magill. Sir Robert Macintosh (1897–1989) introduced a curved laryngoscope blade in 1943; the Macintosh blade remains to this day the most widely used laryngoscope blade for orotracheal intubation.
Between 1945 and 1952, optical engineers built upon the earlier work of Rudolph Schindler (1888–1968), developing the first gastrocamera. In 1964, optical fiber technology was applied to one of these early gastrocameras to produce the first flexible fiberoptic endoscope. Initially used in upper GI endoscopy, this device was first used for laryngoscopy and tracheal intubation by Peter Murphy, an English anesthetist, in 1967. The concept of using a stylet for replacing or exchanging orotracheal tubes was introduced by Finucane and Kupshik in 1978, using a central venous catheter.
By the mid-1980s, the flexible fiberoptic bronchoscope had become an indispensable instrument within the pulmonology and anesthesia communities. The digital revolution of the 21st century has brought newer technology to the art and science of tracheal intubation. Several manufacturers have developed video laryngoscopes which employ digital technology such as the CMOS active pixel sensor (CMOS APS) to generate a view of the glottis so that the trachea may be intubated.
| Biology and health sciences | Treatments | Health |
146444 | https://en.wikipedia.org/wiki/Opabinia | Opabinia | Opabinia regalis is an extinct, stem group arthropod found in the Middle Cambrian Burgess Shale Lagerstätte (505 million years ago) of British Columbia. Opabinia was a soft-bodied animal, measuring up to 7 cm in body length, and had a segmented trunk with flaps along its sides and a fan-shaped tail. The head showed unusual features: five eyes, a mouth under the head and facing backwards, and a clawed proboscis that most likely passed food to its mouth. Opabinia lived on the seafloor, using the proboscis to seek out small, soft food. Fewer than twenty good specimens have been described; 3 specimens of Opabinia are known from the Greater Phyllopod bed, where they constitute less than 0.1% of the community.
When the first thorough examination of Opabinia in 1975 revealed its unusual features, it was thought to be unrelated to any known phylum, or perhaps a relative of arthropod and annelid ancestors. However, later studies since late 1990s consistently support its affinity as a member of basal arthropods, alongside the closely related radiodonts (Anomalocaris and relatives) and gilled lobopodians (Kerygmachela and Pambdelurion).
In the 1970s, there was an ongoing debate about whether multi-celled animals appeared suddenly during the Early Cambrian, in an event called the Cambrian explosion, or had arisen earlier but without leaving fossils. At first Opabinia was regarded as strong evidence for the "explosive" hypothesis. Later the discovery of a whole series of similar lobopodian animals, some with closer resemblances to arthropods, and the development of the idea of stem groups, suggested that the Early Cambrian was a time of relatively fast evolution, but one that could be understood without assuming any unique evolutionary processes.
History of discovery
In 1911, Charles Doolittle Walcott found in the Burgess Shale nine almost complete fossils of Opabinia regalis and a few of what he classified as Opabinia ? media, and published a description of all of these in 1912. The generic name is derived from Opabin pass between Mount Hungabee and Mount Biddle, southeast of Lake O'Hara, British Columbia, Canada. In 1966–1967, Harry B. Whittington found another good specimen, and in 1975 he published a detailed description based on very thorough dissection of some specimens and photographs of these specimens lit from a variety of angles. Whittington's analysis did not cover Opabinia ? media; Walcott's specimens of this species could not be identified in his collection. In 1960 Russian paleontologists described specimens they found in the Norilsk region of Siberia and labelled Opabinia norilica, but these fossils were poorly preserved, and Whittington did not feel they provided enough information to be classified as members of the genus Opabinia.
Occurrence
All the recognized Opabinia specimens found so far come from the "Phyllopod bed" of the Burgess Shale, in the Canadian Rockies of British Columbia.
In 1997, Briggs and Nedin reported from South Australia Emu Bay Shale a new specimen of Myoscolex that was much better preserved than previous specimens, leading them to conclude that it was a close relative of Opabinia—although this interpretation was later questioned by Dzik, who instead concluded that Myoscolex was an annelid worm.
Morphology
Opabinia looked so strange that the audience at the first presentation of Whittington's analysis laughed. The length of Opabinia regalis from head (excluding proboscis) to tail end ranged between and . One of the most distinctive characters of Opabinia is the hollow proboscis, whose total length was about one-third that of the body, and projected down from under the head. The proboscis was striated like a vacuum cleaner's hose and flexible, and it ended with a claw-like structure whose terminal edges bore 5 spines that projected inwards and forwards. The bilateral symmetry and lateral (instead of vertical as reconstructed by Whittington 1975) arrangement of the claw suggest it represents a pair of fused frontal appendages, comparable to those of radiodonts and gilled lobopodians. The head bore five stalked eyes: two near the front and fairly close to the middle of the head, pointing upwards and forwards; two larger eyes with longer stalks near the rear and outer edges of the head, pointing upwards and sideways; and a single eye between the larger pair of stalked eyes, pointing upwards. It has been assumed that the eyes were all compound, like other arthropods' lateral eyes, but this reconstruction, which is not backed up by any evidence, is "somewhat fanciful". The mouth was under the head, behind the proboscis, and pointed backwards, so that the digestive tract formed a U-bend on its way towards the rear of the animal. The proboscis appears to have been sufficiently long and flexible to reach the mouth.
The main part of the body was typically about wide and had 15 segments, on each of which there were pairs of flaps (lobes) pointing downwards and outwards. The flaps overlapped so that the front of each was covered by the rear edge of the one ahead of it. The body ended with what looked like a single conical segment bearing three pairs of overlapping tail fan blades that pointed up and out, forming a tail like a V-shaped double fan.
Interpretations of other features of Opabinia fossils differ. Since the animals did not have mineralized armor nor even tough organic exoskeletons like those of other arthropods, their bodies were flattened as they were buried and fossilized, and smaller or internal features appear as markings within the outlines of the fossils.
Whittington (1975) interpreted the gills as paired extensions attached dorsally to the bases of all but the first flaps on each side, and thought that these gills were flat underneath, had overlapping layers on top. Bergström (1986) revealed the "overlapping layers" were rows of individual blades, interpreted the flaps as part of dorsal coverings (tergite) over the upper surface of the body, with blades attached underneath each of them. Budd (1996) thought the gill blades attached along the front edges on the dorsal side of all except the first flaps. He also found marks inside the flaps' front edges that he interpreted as internal channels connecting the gills to the interior of the body, much as Whittington interpreted the mark along the proboscis as an internal channel. Zhang and Briggs (2007) however, interpreted all flaps have posterior spacing where the gill blades attached. Budd and Daley (2011) reject the reconstruction by Zhang & Briggs, showing the flaps have complete posterior edges as in previous reconstructions. They mostly follow the reconstruction by Budd (1996) with modifications on some details (e.g. the first flap pair also have gills; the attachment point of gill blades located more posteriorly than previously thought).
Whittington (1975) found evidence of near-triangular features along the body, and concluded that they were internal structures, most likely sideways extensions of the gut (diverticula). Chen et al. (1994) interpreted them as contained within the lobes along the sides. Budd (1996) thought the "triangles" were too wide to fit within Opabinias slender body, and that cross-section views showed they were attached separately from and lower than the lobes, and extended below the body. He later found specimens that appeared to preserve the legs' exterior cuticle. He therefore interpreted the "triangles" as short, fleshy, conical legs (lobopods). He also found small mineralized patches at the tips of some, and interpreted these as claws. Under this reconstruction, the gill-bearing flap and lobopod were homologized to the outer gill branch and inner leg branch of arthropod biramous limbs seen in Marrella, trilobites, and crustaceans. Zhang and Briggs (2007) analyzed the chemical composition of the "triangles", and concluded that they had the same composition as the gut, and therefore agreed with Whittington that they were part of the digestive system. Instead they regarded Opabinias lobe+gill arrangement as an early form of the arthropod limbs before it split into a biramous structure. However, this similar chemical composition is not only associated with the digestive tract; Budd and Daley (2011) suggest that it represents mineralization forming within fluid-filled cavities within the body, which is consistent with hollow lobopods as seen in unequivocal lobopodian fossils. They also clarify that the gut diverticula of Opabinia are series of circular gut glands individualized from the "triangles". While they agreed on the absence of terminal claws, the presence of lobopods in Opabinia remain as a plausible interpretation.
Lifestyle
The way in which the Burgess Shale animals were buried, by a mudslide or a sediment-laden current that acted as a sandstorm, suggests they lived on the surface of the seafloor. Opabinia probably used its proboscis to search the sediment for food particles and pass them to its mouth. Since there is no sign of anything that might function as jaws, its food was presumably small and soft. The paired gut diverticula may increase the efficiency of food digestion and intake of nutrition. Whittington (1975) believing that Opabinia had no legs, thought that it crawled on its lobes and that it could also have swum slowly by flapping the lobes, especially if it timed the movements to create a wave with the metachronal movement of its lobes. On the other hand, he thought the body was not flexible enough to allow fish-like undulations of the whole body.
Classification
Considering how paleontologists' reconstructions of Opabinia differ, it is not surprising that the animal's classification was highly debated during the 20th century. Charles Doolittle Walcott, the original describer, considered it to be an anostracan crustacean in 1912. The idea was followed by G. Evelyn Hutchinson in 1930, providing the first reconstruction of Opabinia as an anostracan swimming upside down. Alberto Simonetta provided a new reconstruction of Opabinia in 1970 very different to those of Hutchinson's, with lots of arthropod features (e.g. ,dorsal exoskeleton and jointed limbs) which are reminiscent of Yohoia and Leanchoilia. Leif Størmer, following earlier work by Percy Raymond, thought that Opabinia belonged to the so-called "trilobitoids" (trilobites and similar taxa). After his thorough analysis Harry B. Whittington concluded that Opabinia was not arthropod in 1975, as he found no evidence for arthropodan jointed limbs, and that nothing like the flexible, probably fluid-filled, proboscis was known in arthropods. Although he left Opabinia'''s classification above the family level open, the annulated but not articulated body and the unusual lateral flaps with gills persuaded him that it may have been a representative of the ancestral stock from the origin of annelids and arthropods, two distinct animal phyla (Lophotrochozoan and Ecdysozoan, respectively) which were still thought to be close relatives (united under Articulata) at that time.
In 1985, Derek Briggs and Whittington published a major redescription of Anomalocaris, also from the Burgess Shale. Soon after that, Swedish palaeontologist Jan Bergström, noting in 1986 the similarity of Anomalocaris and Opabinia, suggested that the two animals were related, as they shared numerous features (e.g., lateral flaps, gill blades, stalked eyes, and specialized frontal appendages). He classified them as primitive arthropods, although he considered that arthropods are not a single phylum.
In 1996, Graham Budd found what he considered evidence of short, un-jointed legs in Opabinia. His examination of the gilled lobopodian Kerygmachela from the Sirius Passet lagerstätte, about and over 10M years older than the Burgess Shale, convinced him that this specimen had similar legs. He considered the legs of these two genera very similar to those of the Burgess Shale lobopodian Aysheaia and the modern onychophorans (velvet worms), which are regarded as the bearers of numerous ancestral traits shared by the ancestors of arthropods. After examining several sets of features shared by these and similar lobopodians he drew up a "broad-scale reconstruction of the arthropod stem-group", i.e., of arthropods and what he considered to be their evolutionary basal members. One striking feature of this family tree is that modern tardigrades (water bears) may be Opabinia's closest living evolutionary relatives. On the other hand, Hou et al. (1995, 2006) suggested Opabinia is a member of unusual cycloneuralian worms with convergent arthropod features.
Although Zhang and Briggs (2007) disagreed with Budd's diagnosis that Opabinias "triangles" were legs, the resemblance they saw between Opabinias lobe+gill arrangement and arthropods' biramous limbs led them to conclude that Opabinia was very closely related to arthropods. In fact they presented a family tree very similar to Budd's except that theirs did not mention tardigrades. Regardless of the different morphological interpretations, all major restudies since 1980s similarly concluded that the resemblance between Opabinia and arthropods (e.g., stalked eyes, dorsal segmentation, posterior mouth, fused appendages, gill-like limb branches) are taxonomically significant.
Since the 2010s, the suggested close relationship between Opabinia and tardigrades/cycloneuralians is no longer supported, while the affinity of Opabinia as a stem-group arthropod alongside Radiodonta (a clade that includes Anomalocaris and its relatives) and gilled lobopodians is widely accepted, as consistently shown by multiple phylogenetic analyses, as well as new discoveries such as the presence of arthropod-like gut glands and the intermediate taxon Kylinxia.
In 2022, Paleontologists described a similar looking animal which was discovered in Cambrian-aged rocks of Utah. The fossil was named Utaurora comosa, and was found within the Wheeler Shale. The stem-arthropod was actually first described in 2008, but at the time it was originally considered a specimen of Anomalocaris. This discovery could suggest there were other animals that looked like Opabinia, and its family may have been more diverse.
Theoretical significance
Opabinia made it clear how little was known about soft-bodied animals, which do not usually leave fossils. When Whittington described it in the mid-1970s, there was already a vigorous debate about the early evolution of animals. Preston Cloud argued in 1948 and 1968 that the process was "explosive", and in the early 1970s Niles Eldredge and Stephen Jay Gould developed their theory of punctuated equilibrium, which views evolution as long intervals of near-stasis "punctuated" by short periods of rapid change. On the other hand, around the same time Wyatt Durham and Martin Glaessner both argued that the animal kingdom had a long Proterozoic history that was hidden by the lack of fossils. and Whittington (1975) concluded that Opabinia, and other taxa such as Marrella and Yohoia, cannot be accommodated in modern groups. This was one of the primary reasons why Gould in his book on the Burgess Shale, Wonderful Life, considered that Early Cambrian life was much more disparate and "experimental" than any later set of animals and that the Cambrian explosion was a truly dramatic event, possibly driven by unusual evolutionary mechanisms. He regarded Opabinia as so important to understanding this phenomenon that he wanted to call his book Homage to Opabinia.
However, other discoveries and analyses soon followed, revealing similar-looking animals such as Anomalocaris from the Burgess Shale and Kerygmachela from Sirius Passet. Another Burgess Shale animal, Aysheaia, was considered very similar to modern Onychophora, which are regarded as close relatives of arthropods. Paleontologists defined a group called lobopodians to include fossil panarthropods that are thought to be close relatives of onychophorans, tardigrades and arthropods but lack jointed limbs. This group was later widely accepted as a paraphyletic grade that led to the origin of extant panarthropod phyla.
While this discussion about specific fossils such as Opabinia and Anomalocaris was going on in the late 20th century, the concept of stem groups was introduced to cover evolutionary "aunts" and "cousins". A crown group is a group of closely related living animals plus their last common ancestor plus all its descendants. A stem group contains offshoots from members of the lineage earlier than the last common ancestor of the crown group; it is a relative concept, for example tardigrades are living animals that form a crown group in their own right, but Budd (1996) regarded them also as being a stem group relative to the arthropods. Viewing strange-looking organisms like Opabinia'' in this way makes it possible to see that, while the Cambrian explosion was unusual, it can be understood in terms of normal evolutionary processes.
| Biology and health sciences | Fossil arthropods | Animals |
146539 | https://en.wikipedia.org/wiki/Senescence | Senescence | Senescence () or biological aging is the gradual deterioration of functional characteristics in living organisms. Whole organism senescence involves an increase in death rates or a decrease in fecundity with increasing age, at least in the later part of an organism's life cycle. However, the resulting effects of senescence can be delayed. The 1934 discovery that calorie restriction can extend lifespans by 50% in rats, the existence of species having negligible senescence, and the existence of potentially immortal organisms such as members of the genus Hydra have motivated research into delaying senescence and thus age-related diseases. Rare human mutations can cause accelerated aging diseases.
Environmental factors may affect aging – for example, overexposure to ultraviolet radiation accelerates skin aging. Different parts of the body may age at different rates and distinctly, including the brain, the cardiovascular system, and muscle. Similarly, functions may distinctly decline with aging, including movement control and memory. Two organisms of the same species can also age at different rates, making biological aging and chronological aging distinct concepts.
Definition and characteristics
Organismal senescence is the aging of whole organisms. Actuarial senescence can be defined as an increase in mortality or a decrease in fecundity with age. The Gompertz–Makeham law of mortality says that the age-dependent component of the mortality rate increases exponentially with age.
Aging is characterized by the declining ability to respond to stress, increased homeostatic imbalance, and increased risk of aging-associated diseases including cancer and heart disease. Aging has been defined as "a progressive deterioration of physiological function, an intrinsic age-related process of loss of viability and increase in vulnerability."
In 2013, a group of scientists defined nine hallmarks of aging that are common between organisms with emphasis on mammals:
genomic instability,
telomere attrition,
epigenetic alterations,
loss of proteostasis,
deregulated nutrient sensing,
mitochondrial dysfunction,
cellular senescence,
stem cell exhaustion,
altered intercellular communication
In a decadal update, three hallmarks have been added, totaling 12 proposed hallmarks:
disabled macroautophagy
chronic inflammation
dysbiosis
The environment induces damage at various levels, e.g. damage to DNA, and damage to tissues and cells by oxygen radicals (widely known as free radicals), and some of this damage is not repaired and thus accumulates with time. Cloning from somatic cells rather than germ cells may begin life with a higher initial load of damage. Dolly the sheep died young from a contagious lung disease, but data on an entire population of cloned individuals would be necessary to measure mortality rates and quantify aging.
The evolutionary theorist George Williams wrote, "It is remarkable that after a seemingly miraculous feat of morphogenesis, a complex metazoan should be unable to perform the much simpler task of merely maintaining what is already formed."
Variation among species
Different speeds with which mortality increases with age correspond to different maximum life span among species. For example, a mouse is elderly at 3 years, a human is elderly at 80 years, and ginkgo trees show little effect of age even at 667 years.
Almost all organisms senesce, including bacteria which have asymmetries between "mother" and "daughter" cells upon cell division, with the mother cell experiencing aging, while the daughter is rejuvenated. There is negligible senescence in some groups, such as the genus Hydra. Planarian flatworms have "apparently limitless telomere regenerative capacity fueled by a population of highly proliferative adult stem cells." These planarians are not biologically immortal, but rather their death rate slowly increases with age. Organisms that are thought to be biologically immortal would, in one instance, be Turritopsis dohrnii, also known as the "immortal jellyfish", due to its ability to revert to its youth when it undergoes stress during adulthood. The reproductive system is observed to remain intact, and even the gonads of Turritopsis dohrnii are existing.
Some species exhibit "negative senescence", in which reproduction capability increases or is stable, and mortality falls with age, resulting from the advantages of increased body size during aging.
Theories of aging
More than 300 different theories have been posited to explain the nature (mechanisms) and causes (reasons for natural emergence or factors) of aging. Good theories would both explain past observations and predict the results of future experiments. Some of the theories may complement each other, overlap, contradict, or may not preclude various other theories.
Theories of aging fall into two broad categories, evolutionary theories of aging and mechanistic theories of aging. Evolutionary theories of aging primarily explain why aging happens, but do not concern themselves with the molecular mechanism(s) that drive the process. All evolutionary theories of aging rest on the basic mechanisms that the force of natural selection declines with age. Mechanistic theories of aging can be divided into theories that propose aging is programmed, and damage accumulation theories, i.e. those that propose aging to be caused by specific molecular changes occurring over time.
Evolutionary aging theories
Antagonistic pleiotropy
One theory was proposed by George C. Williams and involves antagonistic pleiotropy. A single gene may affect multiple traits. Some traits that increase fitness early in life may also have negative effects later in life. But, because many more individuals are alive at young ages than at old ages, even small positive effects early can be strongly selected for, and large negative effects later may be very weakly selected against. Williams suggested the following example: Perhaps a gene codes for calcium deposition in bones, which promotes juvenile survival and will therefore be favored by natural selection; however, this same gene promotes calcium deposition in the arteries, causing negative atherosclerotic effects in old age. Thus, harmful biological changes in old age may result from selection for pleiotropic genes that are beneficial early in life but harmful later on. In this case, selection pressure is relatively high when Fisher's reproductive value is high and relatively low when Fisher's reproductive value is low.
Cancer versus cellular senescence tradeoff theory of aging
Senescent cells within a multicellular organism can be purged by competition between cells, but this increases the risk of cancer. This leads to an inescapable dilemma between two possibilities—the accumulation of physiologically useless senescent cells, and cancer—both of which lead to increasing rates of mortality with age.
Disposable soma
The disposable soma theory of aging was proposed by Thomas Kirkwood in 1977. The theory suggests that aging occurs due to a strategy in which an individual only invests in maintenance of the soma for as long as it has a realistic chance of survival. A species that uses resources more efficiently will live longer, and therefore be able to pass on genetic information to the next generation. The demands of reproduction are high, so less effort is invested in repair and maintenance of somatic cells, compared to germline cells, in order to focus on reproduction and species survival.
Programmed aging theories
Programmed theories of aging posit that aging is adaptive, normally invoking selection for evolvability or group selection.
The reproductive-cell cycle theory suggests that aging is regulated by changes in hormonal signaling over the lifespan.
Damage accumulation theories
The free radical theory of aging
One of the most prominent theories of aging was first proposed by Harman in 1956. It posits that free radicals produced by dissolved oxygen, radiation, cellular respiration and other sources cause damage to the molecular machines in the cell and gradually wear them down. This is also known as oxidative stress.
There is substantial evidence to back up this theory. Old animals have larger amounts of oxidized proteins, DNA and lipids than their younger counterparts.
Chemical damage
One of the earliest aging theories was the Rate of Living Hypothesis described by Raymond Pearl in 1928 (based on earlier work by Max Rubner), which states that fast basal metabolic rate corresponds to short maximum life span.
While there may be some validity to the idea that for various types of specific damage detailed below that are by-products of metabolism, all other things being equal, a fast metabolism may reduce lifespan, in general this theory does not adequately explain the differences in lifespan either within, or between, species. Calorically restricted animals process as much, or more, calories per gram of body mass, as their ad libitum fed counterparts, yet exhibit substantially longer lifespans. Similarly, metabolic rate is a poor predictor of lifespan for birds, bats and other species that, it is presumed, have reduced mortality from predation, and therefore have evolved long lifespans even in the presence of very high metabolic rates. In a 2007 analysis it was shown that, when modern statistical methods for correcting for the effects of body size and phylogeny are employed, metabolic rate does not correlate with longevity in mammals or birds.
With respect to specific types of chemical damage caused by metabolism, it is suggested that damage to long-lived biopolymers, such as structural proteins or DNA, caused by ubiquitous chemical agents in the body such as oxygen and sugars, are in part responsible for aging. The damage can include breakage of biopolymer chains, cross-linking of biopolymers, or chemical attachment of unnatural substituents (haptens) to biopolymers.
Under normal aerobic conditions, approximately 4% of the oxygen metabolized by mitochondria is converted to superoxide ion, which can subsequently be converted to hydrogen peroxide, hydroxyl radical and eventually other reactive species including other peroxides and singlet oxygen, which can, in turn, generate free radicals capable of damaging structural proteins and DNA. Certain metal ions found in the body, such as copper and iron, may participate in the process. (In Wilson's disease, a hereditary defect that causes the body to retain copper, some of the symptoms resemble accelerated senescence.) These processes termed oxidative stress are linked to the potential benefits of dietary polyphenol antioxidants, for example in coffee, and tea. However their typically positive effects on lifespans when consumption is moderate have also been explained by effects on autophagy, glucose metabolism and AMPK.
Sugars such as glucose and fructose can react with certain amino acids such as lysine and arginine and certain DNA bases such as guanine to produce sugar adducts, in a process called glycation. These adducts can further rearrange to form reactive species, which can then cross-link the structural proteins or DNA to similar biopolymers or other biomolecules such as non-structural proteins. People with diabetes, who have elevated blood sugar, develop senescence-associated disorders much earlier than the general population, but can delay such disorders by rigorous control of their blood sugar levels. There is evidence that sugar damage is linked to oxidant damage in a process termed glycoxidation.
Free radicals can damage proteins, lipids or DNA. Glycation mainly damages proteins. Damaged proteins and lipids accumulate in lysosomes as lipofuscin. Chemical damage to structural proteins can lead to loss of function; for example, damage to collagen of blood vessel walls can lead to vessel-wall stiffness and, thus, hypertension, and vessel wall thickening and reactive tissue formation (atherosclerosis); similar processes in the kidney can lead to kidney failure. Damage to enzymes reduces cellular functionality. Lipid peroxidation of the inner mitochondrial membrane reduces the electric potential and the ability to generate energy. It is probably no accident that nearly all of the so-called "accelerated aging diseases" are due to defective DNA repair enzymes.
It is believed that the impact of alcohol on aging can be partly explained by alcohol's activation of the HPA axis, which stimulates glucocorticoid secretion, long-term exposure to which produces symptoms of aging.
DNA damage
DNA damage was proposed in a 2021 review to be the underlying cause of aging because of the mechanistic link of DNA damage to nearly every aspect of the aging phenotype. Slower rate of accumulation of DNA damage as measured by the DNA damage marker gamma H2AX in leukocytes was found to correlate with longer lifespans in comparisons of dolphins, goats, reindeer, American flamingos and griffon vultures. DNA damage-induced epigenetic alterations, such as DNA methylation and many histone modifications, appear to be of particular importance to the aging process. Evidence for the theory that DNA damage is the fundamental cause of aging was first reviewed in 1981.
Mutation accumulation
Natural selection can support lethal and harmful alleles, if their effects are felt after reproduction. The geneticist J. B. S. Haldane wondered why the dominant mutation that causes Huntington's disease remained in the population, and why natural selection had not eliminated it. The onset of this neurological disease is (on average) at age 45 and is invariably fatal within 10–20 years. Haldane assumed that, in human prehistory, few survived until age 45. Since few were alive at older ages and their contribution to the next generation was therefore small relative to the large cohorts of younger age groups, the force of selection against such late-acting deleterious mutations was correspondingly small. Therefore, a genetic load of late-acting deleterious mutations could be substantial at mutation–selection balance. This concept came to be known as the selection shadow.
Peter Medawar formalised this observation in his mutation accumulation theory of aging. "The force of natural selection weakens with increasing age—even in a theoretically immortal population, provided only that it is exposed to real hazards of mortality. If a genetic disaster... happens late enough in individual life, its consequences may be completely unimportant". Age-independent hazards such as predation, disease, and accidents, called 'extrinsic mortality', mean that even a population with negligible senescence will have fewer individuals alive in older age groups.
Other damage
A study concluded that retroviruses in the human genomes can become awakened from dormant states and contribute to aging which can be blocked by neutralizing antibodies, alleviating "cellular senescence and tissue degeneration and, to some extent, organismal aging".
Stem cell theories of aging
Hematopoietic stem cell aging
Hematopoietic stem cell diversity aging
Hematopoietic mosaic loss of chromosome Y
Biomarkers of aging
If different individuals age at different rates, then fecundity, mortality, and functional capacity might be better predicted by biomarkers than by chronological age. However, graying of hair, face aging, skin wrinkles and other common changes seen with aging are not better indicators of future functionality than chronological age. Biogerontologists have continued efforts to find and validate biomarkers of aging, but success thus far has been limited.
Levels of CD4 and CD8 memory T cells and naive T cells have been used to give good predictions of the expected lifespan of middle-aged mice.
Aging clocks
There is interest in an epigenetic clock as a biomarker of aging, based on its ability to predict human chronological age. Basic blood biochemistry and cell counts can also be used to accurately predict the chronological age. It is also possible to predict the human chronological age using transcriptomic aging clocks.
There is research and development of further biomarkers, detection systems and software systems to measure biological age of different tissues or systems or overall. For example, a deep learning (DL) software using anatomic magnetic resonance images estimated brain age with relatively high accuracy, including detecting early signs of Alzheimer's disease and varying neuroanatomical patterns of neurological aging, and a DL tool was reported as to calculate a person's inflammatory age based on patterns of systemic age-related inflammation.
Aging clocks have been used to evaluate impacts of interventions on humans, including combination therapies. Exmploying aging clocks to identify and evaluate longevity interventions represents a fundamental goal in aging biology research. However, achieving this goal requires overcoming numerous challenges and implementing additional validation steps.
Genetic determinants of aging
A number of genetic components of aging have been identified using model organisms, ranging from the simple budding yeast Saccharomyces cerevisiae to worms such as Caenorhabditis elegans and fruit flies (Drosophila melanogaster). Study of these organisms has revealed the presence of at least two conserved aging pathways.
Gene expression is imperfectly controlled, and it is possible that random fluctuations in the expression levels of many genes contribute to the aging process as suggested by a study of such genes in yeast. Individual cells, which are genetically identical, nonetheless can have substantially different responses to outside stimuli, and markedly different lifespans, indicating the epigenetic factors play an important role in gene expression and aging as well as genetic factors. There is research into epigenetics of aging.
The ability to repair DNA double-strand breaks declines with aging in mice and humans.
A set of rare hereditary (genetics) disorders, each called progeria, has been known for some time. Sufferers exhibit symptoms resembling accelerated aging, including wrinkled skin. The cause of Hutchinson–Gilford progeria syndrome was reported in the journal Nature in May 2003.
This report suggests that DNA damage, not oxidative stress, is the cause of this form of accelerated aging.
A study indicates that aging may shift activity toward short genes or shorter transcript length and that this can be countered by interventions.
Healthspans and aging in society
Healthspan can broadly be defined as the period of one's life that one is healthy, such as free of significant diseases or declines of capacities (e.g. of senses, muscle, endurance and cognition).
Biological aging or the LHG comes with a great cost burden to society, including potentially rising health care costs (also depending on types and costs of treatments). This, along with global quality of life or wellbeing, highlight the importance of extending healthspans.
Many measures that may extend lifespans may simultaneously also extend healthspans, albeit that is not necessarily the case, indicating that "lifespan can no longer be the sole parameter of interest" in related research. While recent life expectancy increases were not followed by "parallel" healthspan expansion, awareness of the concept and issues of healthspan lags as of 2017. Scientists have noted that "[c]hronic diseases of aging are increasing and are inflicting untold costs on human quality of life".
Interventions
| Biology and health sciences | Fields of medicine | Health |
146635 | https://en.wikipedia.org/wiki/Steamship | Steamship | A steamship, often referred to as a steamer, is a type of steam-powered vessel, typically ocean-faring and seaworthy, that is propelled by one or more steam engines that typically move (turn) propellers or paddlewheels. The first steamships came into practical usage during the early 19th century; however, there were exceptions that came before. Steamships usually use the prefix designations of "PS" for paddle steamer or "SS" for screw steamer (using a propeller or screw). As paddle steamers became less common, "SS" is incorrectly assumed by many to stand for "steamship". Ships powered by internal combustion engines use a prefix such as "MV" for motor vessel, so it is not correct to use "SS" for most modern vessels.
As steamships were less dependent on wind patterns, new trade routes opened up. The steamship has been described as a "major driver of the first wave of trade globalization (1870–1913)" and contributor to "an increase in international trade that was unprecedented in human history".
History
Steamships were preceded by smaller vessels, called steamboats, conceived in the first half of the 18th century by Denis Papin, with the first working steamboat and paddle steamer, the Pyroscaphe, from 1783. Once the technology of steam was mastered at this level, steam engines were mounted on larger, and eventually, ocean-going vessels. Becoming reliable, and propelled by screw rather than paddlewheels, the technology changed the design of ships for faster, more economic propulsion.
Paddlewheels as the main motive source became standard on these early vessels. It was an effective means of propulsion under ideal conditions but otherwise had serious drawbacks. The paddle-wheel performed best when it operated at a certain depth, however when the depth of the ship changed from added weight it further submerged the paddle wheel causing a substantial decrease in performance.
Within a few decades of the development of the river and canal steamboat, the first steamships began to cross the Atlantic Ocean. The first sea-going steamboat was Richard Wright's first steamboat Experiment, an ex-French lugger; she steamed from Leeds to Yarmouth in July 1813.
The first iron steamship to go to sea was the 116-ton Aaron Manby, built in 1821 by Aaron Manby at the Horseley Ironworks, and became the first iron-built vessel to put to sea when she crossed the English Channel in 1822, arriving in Paris on 22 June. She carried passengers and freight to Paris in 1822 at an average speed of 8 knots (9 mph, 14 km/h).
The American ship first crossed the Atlantic Ocean arriving in Liverpool, England, on June 20, 1819, although most of the voyage was actually made under sail. The first ship to make the transatlantic trip substantially under steam power may have been the British-built Dutch-owned Curaçao, a wooden 438-ton vessel built in Dover and powered by two 50 hp engines, which crossed from Hellevoetsluis, near Rotterdam on 26 April 1827 to Paramaribo, Surinam on 24 May, spending 11 days under steam on the way out and more on the return. Another claimant is the Canadian ship in 1833.
The first steamship purpose-built for regularly scheduled trans-Atlantic crossings was the British side-wheel paddle steamer built by Isambard Kingdom Brunel in 1838, which inaugurated the era of the trans-Atlantic ocean liner.
, built in Britain in 1839 by Francis Pettit Smith, was the world's first screw propeller-driven steamship for open water seagoing. She had considerable influence on ship development, encouraging the adoption of screw propulsion by the Royal Navy, in addition to her influence on commercial vessels. The first screw-driven propeller steamship introduced in America was on a ship built by Thomas Clyde in 1844 and many more ships and routes followed.
Screw-propeller steamers
The key innovation that made ocean-going steamers viable was the change from the paddle-wheel to the screw-propeller as the mechanism of propulsion. These steamships quickly became more popular, because the propeller's efficiency was consistent regardless of the depth at which it operated. Being smaller in size and mass and being completely submerged, it was also far less prone to damage.
James Watt of Scotland is widely given credit for applying the first screw propeller to an engine at his Birmingham works, an early steam engine, beginning the use of a hydrodynamic screw for propulsion.
The development of screw propulsion relied on the following technological innovations.
Steam engines had to be designed with the power delivered at the bottom of the machinery, to give direct drive to the propeller shaft. A paddle steamer's engines drive a shaft that is positioned above the waterline, with the cylinders positioned below the shaft. used chain drive to transmit power from a paddler's engine to the propeller shaft – the result of a late design change to propeller propulsion.
An effective stern tube and associated bearings were required. The stern tube contains the propeller shaft where it passes through the hull structure. It should provide an unrestricted delivery of power by the propeller shaft. The combination of hull and stern tube must avoid any flexing that will bend the shaft or cause uneven wear. The inboard end has a stuffing box that prevents water from entering the hull along the tube. Some early stern tubes were made of brass and operated as a water lubricated bearing along the entire length. In other instances a long bush of soft metal was fitted in the after end of the stern tube. had this arrangement fail on her first transatlantic voyage, with very large amounts of uneven wear. The problem was solved with a lignum vitae water-lubricated bearing, patented in 1858. This became standard practice and is in use today.
Since the motive power of screw propulsion is delivered along the shaft, a thrust bearing is needed to transfer that load to the hull without excessive friction. had a 2 ft diameter gunmetal plate on the forward end of the shaft which bore against a steel plate attached to the engine beds. Water at 200 psi was injected between these two surfaces to lubricate and separate them. This arrangement was not sufficient for higher engine powers and oil lubricated "collar" thrust bearings became standard from the early 1850s. This was superseded at the beginning of the 20th century by floating pad bearing which automatically built up wedges of oil which could withstand bearing pressures of 500 psi or more.
Name prefix
Steam-powered ships were named with a prefix designating their propeller configuration i.e. single, twin, triple-screw. Single-screw Steamship SS, Twin-Screw Steamship TSS, Triple-Screw Steamship TrSS. Steam turbine-driven ships had the prefix TS. In the UK the prefix RMS for Royal Mail Steamship overruled the screw configuration prefix.
First ocean-going steamships
The first steamship credited with crossing the Atlantic Ocean between North America and Europe was the American ship , though she was actually a hybrid between a steamship and a sailing ship, with the first half of the journey making use of the steam engine. Savannah left the port of Savannah, Georgia, US, on 22 May 1819, arriving in Liverpool, England, on 20 June 1819; her steam engine having been in use for part of the time on 18 days (estimates vary from 8 to 80 hours). A claimant to the title of the first ship to make the transatlantic trip substantially under steam power is the British-built Dutch-owned Curaçao, a wooden 438-ton vessel built in Dover and powered by two 50 hp engines, which crossed from Hellevoetsluis, near Rotterdam on 26 April 1827 to Paramaribo, Surinam on 24 May, spending 11 days under steam on the way out and more on the return. Another claimant is the Canadian ship in 1833.
The British side-wheel paddle steamer was the first steamship purpose-built for regularly scheduled trans-Atlantic crossings, starting in 1838. In 1836 Isambard Kingdom Brunel and a group of Bristol investors formed the Great Western Steamship Company to build a line of steamships for the Bristol-New York route. The idea of regular scheduled transatlantic service was under discussion by several groups and the rival British and American Steam Navigation Company was established at the same time. Great Western's design sparked controversy from critics that contended that she was too big. The principle that Brunel understood was that the carrying capacity of a hull increases as the cube of its dimensions, while water resistance only increases as the square of its dimensions. This meant that large ships were more fuel efficient, something very important for long voyages across the Atlantic.
Great Western was an iron-strapped, wooden, side-wheel paddle steamer, with four masts to hoist the auxiliary sails. The sails were not just to provide auxiliary propulsion, but also were used in rough seas to keep the ship on an even keel and ensure that both paddle wheels remained in the water, driving the ship in a straight line. The hull was built of oak by traditional methods. She was the largest steamship for one year, until the British and American's British Queen went into service. Built at the shipyard of Patterson & Mercer in Bristol, Great Western was launched on 19 July 1837 and then sailed to London, where she was fitted with two side-lever steam engines from the firm of Maudslay, Sons & Field, producing 750 indicated horsepower between them. The ship proved satisfactory in service and initiated the transatlantic route, acting as a model for all following Atlantic paddle-steamers.
The Cunard Line's began her first regular passenger and cargo service by a steamship in 1840, sailing from Liverpool to Boston.
In 1845 the revolutionary , also built by Brunel, became the first iron-hulled screw-driven ship to cross the Atlantic. SS Great Britain was the first ship to combine these two innovations. After the initial success of its first liner, of 1838, the Great Western Steamship Company assembled the same engineering team that had collaborated so successfully before. This time however, Brunel, whose reputation was at its height, came to assert overall control over design of the ship—a state of affairs that would have far-reaching consequences for the company. Construction was carried out in a specially adapted dry dock in Bristol, England.
Brunel was given a chance to inspect John Laird's (English) channel packet ship Rainbow—the largest iron-hulled ship then in service—in 1838, and was soon converted to iron-hulled technology. He scrapped his plans to build a wooden ship and persuaded the company directors to build an iron-hulled ship. Iron's advantages included being much cheaper than wood, not being subject to dry rot or woodworm, and its much greater structural strength. The practical limit on the length of a wooden-hulled ship is about 300 feet, after which hogging—the flexing of the hull as waves pass beneath it—becomes too great. Iron hulls are far less subject to hogging, so that the potential size of an iron-hulled ship is much greater.
In the spring of 1840 Brunel also had the opportunity to inspect , the first screw-propelled steamship, completed only a few months before by F. P. Smith's Propeller Steamship Company. Brunel had been looking into methods of improving the performance of Great Britains paddlewheels, and took an immediate interest in the new technology, and Smith, sensing a prestigious new customer for his own company, agreed to lend Archimedes to Brunel for extended tests. Over several months, Smith and Brunel tested a number of different propellers on Archimedes in order to find the most efficient design, a four-bladed model submitted by Smith. When launched in 1843, Great Britain was by far the largest vessel afloat.
Brunel's last major project, , was built in 1854–1857 with the intent of linking Great Britain with India, via the Cape of Good Hope, without any coaling stops. This ship was arguably more revolutionary than her predecessors. She was one of the first ships to be built with a double hull with watertight compartments and was the first liner to have four funnels. She was the biggest liner throughout the rest of the 19th century with a gross tonnage of almost 20,000 tons and had a passenger-carrying capacity of thousands. The ship was ahead of her time and went through a turbulent history, never being put to her intended use. The first transatlantic steamer built of steel was , built by Allan Line Royal Mail Steamers and entering service in 1879.
The first regular steamship service from the East Coast to the West Coast of the United States began on 28 February 1849, with the arrival of in San Francisco Bay. The California left New York Harbor on 6 October 1848, rounded Cape Horn at the tip of South America, and arrived at San Francisco, California, after a four-month and 21-day journey. The first steamship to operate on the Pacific Ocean was the paddle steamer Beaver, launched in 1836 to service Hudson's Bay Company trading posts between Puget Sound Washington and Alaska.
Long-distance commercial steamships
The most testing route for steam was from Britain or the East Coast of the U.S. to the Far East. The distance from either is roughly the same, between , traveling down the Atlantic, around the southern tip of Africa, and across the Indian Ocean. Before 1866, no steamship could carry enough coal to make this voyage and have enough space left to carry a commercial cargo.
A partial solution to this problem was adopted by the Peninsular and Oriental Steam Navigation Company (P&O), using an overland section between Alexandria and Suez, with connecting steamship routes along the Mediterranean and then through the Red Sea. While this worked for passengers and some high value cargo, sail was still the only solution for virtually all trade between China and Western Europe or East Coast America. Most notable of these cargoes was tea, typically carried in clippers.
Another partial solution was the Steam Auxiliary Ship – a vessel with a steam engine, but also rigged as a sailing vessel. The steam engine would only be used when conditions were unsuitable for sailing – in light or contrary winds. Some of this type (for instance Erl King) were built with propellers that could be lifted clear of the water to reduce drag when under sail power alone. These ships struggled to be successful on the route to China, as the standing rigging required when sailing was a handicap when steaming into a head wind, most notably against the southwest monsoon when returning with a cargo of new tea. Though the auxiliary steamers persisted in competing in far eastern trade for a few years (and it was Erl King that carried the first cargo of tea through the Suez Canal), they soon moved on to other routes.
What was needed was a big improvement in fuel efficiency. While the boilers for steam engines on land were allowed to run at high pressures, the Board of Trade (under the authority of the Merchant Shipping Act 1854) would not allow ships to exceed . Compound engines were a known source of improved efficiency – but generally not used at sea due to the low pressures available. Carnatic (1863), a P&O ship, had a compound engine – and achieved better efficiency than other ships of the time. Her boilers ran at but relied on a substantial amount of superheat.
Alfred Holt, who had entered marine engineering and ship management after an apprenticeship in railway engineering, experimented with boiler pressures of in Cleator. Holt was able to persuade the Board of Trade to allow these boiler pressures and, in partnership with his brother Phillip launched Agamemnon in 1865. Holt had designed a particularly compact compound engine and taken great care with the hull design, producing a light, strong, easily driven hull.
The efficiency of Holt's package of boiler pressure, compound engine and hull design gave a ship that could steam at 10 knots on 20 long tons of coal a day. This fuel consumption was a saving from between 23 and 14 long tons a day, compared to other contemporary steamers. Not only did less coal need to be carried to travel a given distance, but fewer firemen were needed to fuel the boilers, so crew costs and their accommodation space were reduced. Agamemnon was able to sail from London to China with a coaling stop at Mauritius on the outward and return journey, with a time on passage substantially less than the competing sailing vessels. Holt had already ordered two sister ships to Agamemnon by the time she had returned from her first trip to China in 1866, operating these ships in the newly formed Blue Funnel Line. His competitors rapidly copied his ideas for their own new ships.
The opening of the Suez Canal in 1869 gave a distance saving of about on the route from China to London. The canal was not a practical option for sailing vessels, as using a tug was difficult and expensive – so this distance saving was not available to them. Steamships immediately made use of this new waterway and found themselves in high demand in China for the start of the 1870 tea season. The steamships were able to obtain a much higher rate of freight than sailing ships and the insurance premium for the cargo was less. So successful were the steamers using the Suez Canal that, in 1871, 45 were built in Clyde shipyards alone for Far Eastern trade.
Triple expansion engines
Throughout the 1870s, compound-engined steamships and sailing vessels coexisted in an economic equilibrium: the operating costs of steamships were still too high in certain trades, so sail was the only commercial option in many situations. The compound engine, where steam was expanded twice in two separate cylinders, still had inefficiencies. The solution was the triple expansion engine, in which steam was successively expanded in a high pressure, intermediate pressure and a low pressure cylinder.
The theory of this was established in the 1850s by John Elder, but it was clear that triple expansion engines needed steam at, by the standards of the day, very high pressures. The existing boiler technology could not deliver this. Wrought iron could not provide the strength for the higher pressures. Steel became available in larger quantities in the 1870s, but the quality was variable. The overall design of boilers was improved in the early 1860s, with the Scotch-type boilers – but at that date these still ran at the lower pressures that were then current.
The first ship fitted with triple expansion engines was Propontis (launched in 1874). She was fitted with boilers that operated at – but these had technical problems and had to be replaced with ones that ran at . This substantially degraded performance.
There were a few further experiments until went into service on the route from Britain to Australia. Her triple expansion engine was designed by Dr A C Kirk, the engineer who had developed the machinery for Propontis. The difference was the use of two double ended Scotch type steel boilers, running at . These boilers had patent corrugated furnaces that overcame the competing problems of heat transfer and sufficient strength to deal with the boiler pressure. Aberdeen was a marked success, achieving in trials, at 1,800 indicated horsepower, a fuel consumption of of coal per indicated horsepower. This was a reduction in fuel consumption of about 60%, compared to a typical steamer built ten years earlier. In service, this translated into less than 40 tons of coal a day when travelling at . Her maiden outward voyage to Melbourne took 42 days, with one coaling stop, carrying 4,000 tons of cargo.
Other similar ships were rapidly brought into service over the next few years. By 1885 the usual boiler pressure was and virtually all ocean-going steamships being built were ordered with triple expansion engines. Within a few years, new installations were running at . The tramp steamers that operated at the end of the 1880s could sail at with a fuel consumption of of coal per ton mile travelled. This level of efficiency meant that steamships could now operate as the primary method of maritime transport in the vast majority of commercial situations. In 1890, steamers constituted 57% of world's tonnage, and by World War I their share raised to 93%.
Era of the ocean liner
By 1870 a number of inventions such as the screw propeller, the compound engine, and the triple-expansion engine made trans-oceanic shipping on a large scale economically viable. In 1870 the White Star Line’s set a new standard for ocean travel by having its first-class cabins amidships, with the added amenity of large portholes, electricity and running water. The size of ocean liners increased from 1880 to meet the needs of the human migration to the United States and Australia.
and her sister ship were the last two Cunard liners of the period to be fitted with auxiliary sails. Both ships were built by John Elder & Co. of Glasgow, Scotland, in 1884. They were record breakers by the standards of the time, and were the largest liners then in service, plying the Liverpool to New York route.
was the largest steamship in the world when she sank in 1912; a subsequent major sinking of a steamer was that of the , as an act of World War I.
Launched in 1938, was the largest passenger steamship ever built. Launched in 1969, Queen Elizabeth 2 (QE2) was the last passenger steamship to cross the Atlantic Ocean on a scheduled liner voyage before she was converted to diesels in 1986. The last major passenger ship built with steam turbines was the Fairsky, launched in 1984, later Atlantic Star, reportedly sold to Turkish shipbreakers in 2013.
Most luxury yachts at the end of the 19th and early 20th centuries were steam driven (see luxury yacht; also Cox & King yachts). Thomas Assheton Smith was an English aristocrat who forwarded the design of the steam yacht in conjunction with the Scottish marine engineer Robert Napier.
Decline of the steamship
By World War II, steamers still constituted 73% of world's tonnage, and similar percentage remained in early 1950s. The decline of the steamship began soon thereafter. Many had been lost in the war, and marine diesel engines had finally matured as an economical and viable alternative to steam power. The diesel engine had far better thermal efficiency than the reciprocating steam engine, and was far easier to control. Diesel engines also required far less supervision and maintenance than steam engines, and as an internal combustion engine it did not need boilers or a water supply, therefore was more space efficient and cheaper to build.
The Liberty ships were the last major steamship class equipped with reciprocating engines. The last Victory ships had already been equipped with marine diesels, and diesel engines superseded both steamers and windjammers soon after World War Two. Most steamers were used up to their maximum economical life span, and no commercial ocean-going steamers with reciprocating engines have been built since the 1960s.
1970–present
Most steamships today are powered by steam turbines. After the demonstration by British engineer Charles Parsons of his steam turbine-driven yacht, Turbinia, in 1897, the use of steam turbines for propulsion quickly spread. The Cunard RMS Mauretania, built in 1906 was one of the first ocean liners to use the steam turbine (with a late design change shortly before her keel was laid down) and was soon followed by all subsequent liners.
Most larger warships of the world's navies were propelled by steam turbines burning bunker fuel in both World Wars, apart from obsolete ships with reciprocating machines from the turn of the century, and rare cases of usage of diesel engines in larger warships. Steam turbines burning fuel remained in warship construction until the end of the Cold War (eg. Russian aircraft carrier Admiral Kuznetsov), because of needs of high power and speed, although from 1970s they were mostly replaced by gas turbines. Large naval vessels and submarines continue to be operated with steam turbines, using nuclear reactors to boil the water.
NS Savannah, was the first nuclear-powered cargo-passenger ship, and was built in the late 1950s as a demonstration project for the potential use of nuclear energy.
Thousands of Liberty Ships (powered by steam piston engines) and Victory Ships (powered by steam turbine engines) were built in World War II. A few of these survive as floating museums and sail occasionally: , , , , and .
A steam turbine ship can be either direct propulsion (the turbines, equipped with a reduction gear, rotate directly the propellers), or turboelectric (the turbines rotate electric generators, which in turn feed electric motors operating the propellers).
While steam turbine-driven merchant ships such as the Algol-class cargo ships (1972–1973), ALP Pacesetter-class container ships (1973–1974) and very large crude carriers were built until the 1970s, the use of steam for marine propulsion in the commercial market has declined dramatically due to the development of more efficient diesel engines. One notable exception are LNG carriers which use boil-off gas from the cargo tanks as fuel. However, even there the development of dual-fuel engines has pushed steam turbines into a niche market with about 10% market share in newbuildings in 2013. Lately, there has been some development in hybrid power plants where the steam turbine is used together with gas engines. As of August 2017 the newest class of Steam Turbine ships are the Seri Camellia-class LNG carriers built by Hyundai Heavy Industries (HHI) starting in 2016 and comprising five units.
Nuclear powered ships are basically steam turbine vessels. The boiler is heated, not by heat of combustion, but by the heat generated by nuclear reactor. Most atomic-powered ships today are either aircraft carriers or submarines.
| Technology | Maritime transport | null |
146806 | https://en.wikipedia.org/wiki/Recurrence%20relation | Recurrence relation | In mathematics, a recurrence relation is an equation according to which the th term of a sequence of numbers is equal to some combination of the previous terms. Often, only previous terms of the sequence appear in the equation, for a parameter that is independent of ; this number is called the order of the relation. If the values of the first numbers in the sequence have been given, the rest of the sequence can be calculated by repeatedly applying the equation.
In linear recurrences, the th term is equated to a linear function of the previous terms. A famous example is the recurrence for the Fibonacci numbers,
where the order is two and the linear function merely adds the two previous terms. This example is a linear recurrence with constant coefficients, because the coefficients of the linear function (1 and 1) are constants that do not depend on For these recurrences, one can express the general term of the sequence as a closed-form expression of . As well, linear recurrences with polynomial coefficients depending on are also important, because many common elementary functions and special functions have a Taylor series whose coefficients satisfy such a recurrence relation (see holonomic function).
Solving a recurrence relation means obtaining a closed-form solution: a non-recursive function of .
The concept of a recurrence relation can be extended to multidimensional arrays, that is, indexed families that are indexed by tuples of natural numbers.
Definition
A recurrence relation is an equation that expresses each element of a sequence as a function of the preceding ones. More precisely, in the case where only the immediately preceding element is involved, a recurrence relation has the form
where
is a function, where is a set to which the elements of a sequence must belong. For any , this defines a unique sequence with as its first element, called the initial value.
It is easy to modify the definition for getting sequences starting from the term of index 1 or higher.
This defines recurrence relation of first order. A recurrence relation of order has the form
where is a function that involves consecutive elements of the sequence.
In this case, initial values are needed for defining a sequence.
Examples
Factorial
The factorial is defined by the recurrence relation
and the initial condition
This is an example of a linear recurrence with polynomial coefficients of order 1, with the simple polynomial (in )
as its only coefficient.
Logistic map
An example of a recurrence relation is the logistic map defined by
for a given constant The behavior of the sequence depends dramatically on but is stable when the initial condition varies.
Fibonacci numbers
The recurrence of order two satisfied by the Fibonacci numbers is the canonical example of a homogeneous linear recurrence relation with constant coefficients (see below). The Fibonacci sequence is defined using the recurrence
with initial conditions
Explicitly, the recurrence yields the equations
etc.
We obtain the sequence of Fibonacci numbers, which begins
0, 1, 1, 2, 3, 5, 8, 13, 21, 34, 55, 89, ...
The recurrence can be solved by methods described below yielding Binet's formula, which involves powers of the two roots of the characteristic polynomial ; the generating function of the sequence is the rational function
Binomial coefficients
A simple example of a multidimensional recurrence relation is given by the binomial coefficients , which count the ways of selecting elements out of a set of elements.
They can be computed by the recurrence relation
with the base cases . Using this formula to compute the values of all binomial coefficients generates an infinite array called Pascal's triangle. The same values can also be computed directly by a different formula that is not a recurrence, but uses factorials, multiplication and division, not just additions:
The binomial coefficients can also be computed with a uni-dimensional recurrence:
with the initial value (The division is not displayed as a fraction for emphasizing that it must be computed after the multiplication, for not introducing fractional numbers).
This recurrence is widely used in computers because it does not require to build a table as does the bi-dimensional recurrence, and does involve very large integers as does the formula with factorials (if one uses all involved integers are smaller than the final result).
Difference operator and difference equations
The is an operator that maps sequences to sequences, and, more generally, functions to functions. It is commonly denoted and is defined, in functional notation, as
It is thus a special case of finite difference.
When using the index notation for sequences, the definition becomes
The parentheses around and are generally omitted, and must be understood as the term of index in the sequence and not applied to the element
Given sequence the of is
The is
A simple computation shows that
More generally: the th difference is defined recursively as and one has
This relation can be inverted, giving
A of order is an equation that involves the first differences of a sequence or a function, in the same way as a differential equation of order relates the first derivatives of a function.
The two above relations allow transforming a recurrence relation of order into a difference equation of order , and, conversely, a difference equation of order into recurrence relation of order . Each transformation is the inverse of the other, and the sequences that are solution of the difference equation are exactly those that satisfies the recurrence relation.
For example, the difference equation
is equivalent to the recurrence relation
in the sense that the two equations are satisfied by the same sequences.
As it is equivalent for a sequence to satisfy a recurrence relation or to be the solution of a difference equation, the two terms "recurrence relation" and "difference equation" are sometimes used interchangeably. See Rational difference equation and Matrix difference equation for example of uses of "difference equation" instead of "recurrence relation"
Difference equations resemble differential equations, and this resemblance is often used to mimic methods for solving differentiable equations to apply to solving difference equations, and therefore recurrence relations.
Summation equations relate to difference equations as integral equations relate to differential equations. See time scale calculus for a unification of the theory of difference equations with that of differential equations.
From sequences to grids
Single-variable or one-dimensional recurrence relations are about sequences (i.e. functions defined on one-dimensional grids). Multi-variable or n-dimensional recurrence relations are about -dimensional grids. Functions defined on -grids can also be studied with partial difference equations.
Solving
Solving linear recurrence relations with constant coefficients
Solving first-order non-homogeneous recurrence relations with variable coefficients
Moreover, for the general first-order non-homogeneous linear recurrence relation with variable coefficients:
there is also a nice method to solve it:
Let
Then
If we apply the formula to and take the limit , we get the formula for first order linear differential equations with variable coefficients; the sum becomes an integral, and the product becomes the exponential function of an integral.
Solving general homogeneous linear recurrence relations
Many homogeneous linear recurrence relations may be solved by means of the generalized hypergeometric series. Special cases of these lead to recurrence relations for the orthogonal polynomials, and many special functions. For example, the solution to
is given by
the Bessel function, while
is solved by
the confluent hypergeometric series. Sequences which are the solutions of linear difference equations with polynomial coefficients are called P-recursive. For these specific recurrence equations algorithms are known which find polynomial, rational or hypergeometric solutions.
Solving first-order rational difference equations
A first order rational difference equation has the form . Such an equation can be solved by writing as a nonlinear transformation of another variable which itself evolves linearly. Then standard methods can be used to solve the linear difference equation in .
Stability
Stability of linear higher-order recurrences
The linear recurrence of order ,
has the characteristic equation
The recurrence is stable, meaning that the iterates converge asymptotically to a fixed value, if and only if the eigenvalues (i.e., the roots of the characteristic equation), whether real or complex, are all less than unity in absolute value.
Stability of linear first-order matrix recurrences
In the first-order matrix difference equation
with state vector and transition matrix , converges asymptotically to the steady state vector if and only if all eigenvalues of the transition matrix (whether real or complex) have an absolute value which is less than 1.
Stability of nonlinear first-order recurrences
Consider the nonlinear first-order recurrence
This recurrence is locally stable, meaning that it converges to a fixed point from points sufficiently close to , if the slope of in the neighborhood of is smaller than unity in absolute value: that is,
A nonlinear recurrence could have multiple fixed points, in which case some fixed points may be locally stable and others locally unstable; for continuous f two adjacent fixed points cannot both be locally stable.
A nonlinear recurrence relation could also have a cycle of period for . Such a cycle is stable, meaning that it attracts a set of initial conditions of positive measure, if the composite function
with appearing times is locally stable according to the same criterion:
where is any point on the cycle.
In a chaotic recurrence relation, the variable stays in a bounded region but never converges to a fixed point or an attracting cycle; any fixed points or cycles of the equation are unstable. | Mathematics | Sequences | null |
146808 | https://en.wikipedia.org/wiki/Non-lethal%20weapon | Non-lethal weapon | Non-lethal weapons , also called nonlethal weapons, less-lethal weapons, less-than-lethal weapons, non-deadly weapons, compliance weapons, or pain-inducing weapons are weapons intended to be less likely to kill a living target than conventional weapons such as knives and firearms with live ammunition. It is often understood that unintended or incidental casualties are risked wherever force is applied; however, non-lethal weapons minimise the risk of casualties (e.g. serious/permanent injuries or death) as much as possible. Non-lethal weapons are used in policing and combat situations to limit the escalation of conflict where employment of lethal force is prohibited or undesirable, where rules of engagement require minimum casualties, or where policy restricts the use of conventional force. However, these weapons occasionally cause serious injuries or death due to allergic reactions, improper use and/or other factors; for this reason the term "less-lethal" has been preferred by some organizations as it describes the risks of death more accurately than the term "non-lethal", which some have argued is a misnomer.
Non-lethal weapons may be used by conventional military in a range of missions across the force continuum. They may also be used by military police, by United Nations forces, and by occupation forces for peacekeeping and stability operations. Non-lethal weapons may also be used to channelize a battlefield, control the movement of civilian populations, or to limit civilian access to restricted areas (as they were utilized by the USMC's 1st Marine Expeditionary Force in Somalia in 1995). Similar weapons, tactics, techniques and procedures are employed by police forces domestically in riot control, prisoner control, crowd control, refugee control, and self-defense, where the terminology of "less-than-lethal" is often used.
History
Military
In the past, military and police faced with undesirable escalation of conflict had few acceptable options. Military personnel guarding embassies often found themselves restricted to carrying unloaded weapons. National guards or policing forces charged with quelling riots were able to use only batons or similar club-like weapons, or bayonet or sword charges, or fire live ammunition at crowds. In the late 1980s and early 1990s, the Non-lethality Policy Review Group at U.S. Global Strategy Council in Washington and other independent think tanks around the world called for a concerted effort to develop weapons that were more life-conserving, environmentally friendly, and fiscally responsible than weapons available at that time. The U.S. Congress and other governments agreed and began an organized development of non-lethal weapons to provide a range of options between talking and shooting.
Recognizing the need to limit the escalation of force, research and development of a range of non-lethal weapons has since been undertaken internationally by governments and weapons manufacturers to fill the need for such weapons. Some non-lethal weapons may provide more effective riot control than firearms, truncheons or bayonets with less risk of loss of life or serious injury. Before the general availability of early military non-lethal weapons in the mid 1990s, war-fighters had few or no casualty-limiting options for the employment of scalable force and were continually at risk whenever lethal force was prohibited during sensitive missions.
In 2001, the United States Marine Corps revealed its development of a less-than-lethal energy weapon called the Active Denial System, a focused high frequency microwave device said to be capable of heating all living matter in the target area rapidly and continuously for the duration of the beam, causing transient intolerable pain but no lasting damage. The skin temperature of a person subjected to this weapon can jump to approximately in as little as 2 seconds depending on the skin's starting temperature. The system is nonlethal (the penetration of the beam into human skin is only a few millimeters).
In 2004, author Jon Ronson cited an unclassified military report titled "Non-Lethal Weapons: Terms and | Technology | Less-lethal weapons | null |
146839 | https://en.wikipedia.org/wiki/Gram | Gram | The gram (originally gramme; SI unit symbol g) is a unit of mass in the International System of Units (SI) equal to one thousandth of a kilogram.
Originally defined as of 1795 as "the absolute weight of a volume of pure water equal to the cube of the hundredth part of a metre [1 cm3], and at the temperature of melting ice", the defining temperature (≈0 °C) was later changed to 4 °C, the temperature of maximum density of water.
By the late 19th century, there was an effort to make the base unit the kilogram and the gram a derived unit. In 1960, the new International System of Units defined a gram as one one-thousandth of a kilogram (i.e., one gram is ). The kilogram, as of 2019, is defined by the International Bureau of Weights and Measures from the fixed numerical value of the Planck constant ().
Official SI symbol
The only unit symbol for gram that is recognised by the International System of Units (SI) is "g" following the numeric value with a space, as in "640 g" to stand for "640 grams" in the English language. The SI disallows use of abbreviations such as "gr" (which is the symbol for grains), "gm" ("g⋅m" is the SI symbol for gram-metre) or "Gm" (the SI symbol for gigametre).
History
The word gramme was adopted by the French National Convention in its 1795 decree revising the metric system as replacing the gravet (introduced in 1793 simultaneously with a base measure called grave, of which gravet was a subdivision). Its definition remained that of the weight of a cubic centimetre of water.
French gramme was taken from the Late Latin term . This word—ultimately from Greek (grámma), "letter"—had adopted a specialised meaning in Late Antiquity of "one twenty-fourth part of an ounce" (two oboli), corresponding to about 1.14 modern grams. This use of the term is found in the carmen de ponderibus et mensuris ("poem about weights and measures") composed around 400 AD.
There is also evidence that the Greek was used in the same sense at around the same time, in the 4th century, and survived in this sense into Medieval Greek, while the Latin term died out in Medieval Latin and was recovered in Renaissance scholarship.
The gram was the base unit of mass in the 19th-century centimetre–gram–second system of units (CGS). The CGS system coexisted with the metre–kilogram–second system of units (MKS), first proposed in 1901, during much of the 20th century, but the gram was displaced by the kilogram as the base unit for mass when the MKS system was chosen for the SI base units in 1960.
Uses
The gram is the most widely used unit of measurement for non-liquid ingredients in cooking and grocery shopping worldwide. Liquid ingredients are often measured by volume rather than mass.
Many standards and legal requirements for nutrition labels on food products require relative contents to be stated per 100 g of the product, such that the resulting figure can also be read as a percentage.
Conversion factors
1 gram (g) ≈ (gr)
1 grain (gr) ≈
1 avoirdupois ounce (oz) ≈
1 troy ounce (ozt) = 31.1034768 g (exact, by definition)
100 grams (g) ≈ (oz)
1 carat (ct) = 0.2 grams
1 gamma (γ) = 10−6 grams
1 undecimogramme = 1 "eleventh-gram" = 10−11 grams in the historical quadrant–eleventh-gram–second system (QES system) a.k.a. hebdometre–undecimogramme–second system (HUS system)
500 grams (g) = 1 jin in the Chinese units of measurement.
Comparisons
1 gram is roughly equal to the mass of 1 small paper clip or pen cap.
The Japanese 1 yen coin has a mass of 1 gram, lighter than the British penny (3.56 g), the United States penny (2.5 g), the Euro cent (2.30 g), and the Australian 5 cent coin (2.80 g).
| Physical sciences | Mass and weight | Basics and measurement |
146879 | https://en.wikipedia.org/wiki/Hypothermia | Hypothermia | Hypothermia is defined as a body core temperature below in humans. Symptoms depend on the temperature. In mild hypothermia, there is shivering and mental confusion. In moderate hypothermia, shivering stops and confusion increases. In severe hypothermia, there may be hallucinations and paradoxical undressing, in which a person removes their clothing, as well as an increased risk of the heart stopping.
Hypothermia has two main types of causes. It classically occurs from exposure to cold weather and cold water immersion. It may also occur from any condition that decreases heat production or increases heat loss. Commonly, this includes alcohol intoxication but may also include low blood sugar, anorexia and advanced age. Body temperature is usually maintained near a constant level of through thermoregulation. Efforts to increase body temperature involve shivering, increased voluntary activity, and putting on warmer clothing. Hypothermia may be diagnosed based on either a person's symptoms in the presence of risk factors or by measuring a person's core temperature.
The treatment of mild hypothermia involves warm drinks, warm clothing, and voluntary physical activity. In those with moderate hypothermia, heating blankets and warmed intravenous fluids are recommended. People with moderate or severe hypothermia should be moved gently. In severe hypothermia, extracorporeal membrane oxygenation (ECMO) or cardiopulmonary bypass may be useful. In those without a pulse, cardiopulmonary resuscitation (CPR) is indicated along with the above measures. Rewarming is typically continued until a person's temperature is greater than . If there is no improvement at this point or the blood potassium level is greater than 12 millimoles per litre at any time, resuscitation may be discontinued.
Hypothermia is the cause of at least 1,500 deaths a year in the United States. It is more common in older people and males. One of the lowest documented body temperatures from which someone with accidental hypothermia has survived is in a 2-year-old boy from Poland named Adam. Survival after more than six hours of CPR has been described. In individuals for whom ECMO or bypass is used, survival is around 50%. Deaths due to hypothermia have played an important role in many wars.
The term is from Greek ῠ̔πο (ypo), meaning "under", and θέρμη (thérmē), meaning "heat". The opposite of hypothermia is hyperthermia, an increased body temperature due to failed thermoregulation.
Classification
Hypothermia is often defined as any body temperature below . With this method it is divided into degrees of severity based on the core temperature.
Another classification system, the Swiss staging system, divides hypothermia based on the presenting symptoms which is preferred when it is not possible to determine an accurate core temperature.
Other cold-related injuries that can be present either alone or in combination with hypothermia include:
Chilblains: condition caused by repeated exposure of skin to temperatures just above freezing. The cold causes damage to small blood vessels in the skin. This damage is permanent and the redness and itching will return with additional exposure. The redness and itching typically occurs on cheeks, ears, fingers, and toes.
Frostbite: the freezing and destruction of tissue, which happens below the freezing point of water
Frostnip: a superficial cooling of tissues without cellular destruction
Trench foot or immersion foot: a condition caused by repetitive exposure to water at non-freezing temperatures
The normal human body temperature is often stated as . Hyperthermia and fevers are defined as a temperature of greater than .
Signs and symptoms
Signs and symptoms vary depending on the degree of hypothermia, and may be divided by the three stages of severity. People with hypothermia may appear pale and feel cold to touch.
Infants with hypothermia may feel cold when touched, with bright red skin and an unusual lack of energy.
Behavioural changes such as impaired judgement, impaired sense of time and place, unusual aggression and numbness can be observed in individuals with hypothermia, they can also deny their condition and refuse any help. A hypothermic person can be euphoric and hallucinating.
Cold stress refers to a near-normal body temperature with low skin temperature; signs include shivering. Cold stress is caused by cold exposure and it can lead to hypothermia and frostbite if not treated.
Mild
Symptoms of mild hypothermia may be vague, with sympathetic nervous system excitation (shivering, high blood pressure, fast heart rate, fast respiratory rate, and contraction of blood vessels). These are all physiological responses to preserve heat. Increased urine production due to cold, mental confusion, and liver dysfunction may also be present. Hyperglycemia may be present, as glucose consumption by cells and insulin secretion both decrease, and tissue sensitivity to insulin may be blunted. Sympathetic activation also releases glucose from the liver. In many cases, however, especially in people with alcoholic intoxication, hypoglycemia appears to be a more common cause. Hypoglycemia is also found in many people with hypothermia, as hypothermia may be a result of hypoglycemia.
Moderate
As hypothermia progresses, symptoms include: mental status changes such as amnesia, confusion, slurred speech, decreased reflexes, and loss of fine motor skills.
Severe
As the temperature decreases, further physiological systems falter and heart rate, respiratory rate, and blood pressure all decrease. This results in an expected heart rate in the 30s at a temperature of .
There is often cold, inflamed skin, hallucinations, lack of reflexes, fixed dilated pupils, low blood pressure, pulmonary edema, and shivering is often absent. Pulse and respiration rates decrease significantly, but fast heart rates (ventricular tachycardia, atrial fibrillation) can also occur. Atrial fibrillation is not typically a concern in and of itself.
Paradoxical undressing
Twenty to fifty percent of hypothermia deaths are associated with paradoxical undressing. This typically occurs during moderate and severe hypothermia, as the person becomes disoriented, confused, and combative. They may begin discarding their clothing, which, in turn, increases the rate of heat loss.
Rescuers who are trained in mountain survival techniques are taught to expect this; however, people who die from hypothermia in urban environments who are found in an undressed state are sometimes incorrectly assumed to have been subjected to sexual assault.
One explanation for the effect is a cold-induced malfunction of the hypothalamus, the part of the brain that regulates body temperature. Another explanation is that the muscles contracting peripheral blood vessels become exhausted (known as a loss of vasomotor tone) and relax, leading to a sudden surge of blood (and heat) to the extremities, causing the person to feel overheated.
Terminal burrowing
An apparent self-protective behaviour, known as "terminal burrowing", or "hide-and-die syndrome", occurs in the final stages of hypothermia. Those affected will enter small, enclosed spaces, such as underneath beds or behind wardrobes. It is often associated with paradoxical undressing. Researchers in Germany claim this is "obviously an autonomous process of the brain stem, which is triggered in the final state of hypothermia and produces a primitive and burrowing-like behavior of protection, as seen in hibernating mammals". This happens mostly in cases where temperature drops slowly.
Causes
Hypothermia usually occurs from exposure to low temperatures, and is frequently complicated by alcohol consumption. Any condition that decreases heat production, increases heat loss, or impairs thermoregulation, however, may contribute. Thus, hypothermia risk factors include: substance use disorders (including alcohol use disorder), homelessness, any condition that affects judgment (such as hypoglycemia), the extremes of age, poor clothing, chronic medical conditions (such as hypothyroidism and sepsis), and living in a cold environment. Hypothermia occurs frequently in major trauma, and is also observed in severe cases of anorexia nervosa. Hypothermia is also associated with worse outcomes in people with sepsiswhile most people with sepsis develop fevers (elevated body temperature), some develop hypothermia.
In urban areas, hypothermia frequently occurs with chronic cold exposure, such as in cases of homelessness, as well as with immersion accidents involving drugs, alcohol or mental illness. While studies have shown that people experiencing homelessness are at risk of premature death from hypothermia, the true incidence of hypothermia-related deaths in this population is difficult to determine. In more rural environments, the incidence of hypothermia is higher among people with significant comorbidities and less ability to move independently. With rising interest in wilderness exploration, and outdoor and water sports, the incidence of hypothermia secondary to accidental exposure may become more frequent in the general population.
Alcohol
Alcohol consumption increases the risk of hypothermia in two ways: vasodilation and temperature controlling systems in the brain. Vasodilation increases blood flow to the skin, resulting in heat being lost to the environment. This produces the effect of feeling warm, when one is actually losing heat. Alcohol also affects the temperature-regulating system in the brain, decreasing the body's ability to shiver and use energy that would normally aid the body in generating heat. The overall effects of alcohol lead to a decrease in body temperature and a decreased ability to generate body heat in response to cold environments. Alcohol is a common risk factor for death due to hypothermia. Between 33% and 73% of hypothermia cases are complicated by alcohol.
Water immersion
Hypothermia continues to be a major limitation to swimming or diving in cold water. The reduction in finger dexterity due to pain or numbness decreases general safety and work capacity, which consequently increases the risk of other injuries.
Other factors predisposing to immersion hypothermia include dehydration, inadequate rewarming between repetitive dives, starting a dive while wearing cold, wet dry suit undergarments, sweating with work, inadequate thermal insulation, and poor physical conditioning.
Heat is lost much more quickly in water than in air. Thus, water temperatures that would be quite reasonable as outdoor air temperatures can lead to hypothermia in survivors, although this is not usually the direct clinical cause of death for those who are not rescued. A water temperature of can lead to death in as little as one hour, and water temperatures near freezing can cause death in as little as 15 minutes. During the sinking of the Titanic, most people who entered the water died in 15–30 minutes.
The actual cause of death in cold water is usually the bodily reactions to heat loss and to freezing water, rather than hypothermia (loss of core temperature) itself. For example, plunged into freezing seas, around 20% of victims die within two minutes from cold shock (uncontrolled rapid breathing, and gasping, causing water inhalation, massive increase in blood pressure and cardiac strain leading to cardiac arrest, and panic); another 50% die within 15–30 minutes from cold incapacitation: inability to use or control limbs and hands for swimming or gripping, as the body "protectively" shuts down the peripheral muscles of the limbs to protect its core. Exhaustion and unconsciousness cause drowning, claiming the rest within a similar time.
Pathophysiology
Heat is primarily generated in muscle tissue, including the heart, and in the liver, while it is lost through the skin (90%) and lungs (10%). Heat production may be increased (to over 1200 W in trained endurance athletes) through muscle contractions (i.e. exercise and shivering). The rate of heat loss is determined, as with any object, by convection, conduction, and radiation. The rates of these can be affected by body mass index, body surface area to volume ratios, clothing and other environmental conditions.
Many changes to physiology occur as body temperatures decrease. These occur in the cardiovascular system leading to the Osborn J wave and other dysrhythmias, decreased central nervous system electrical activity, cold diuresis, and non-cardiogenic pulmonary edema.
Research has shown that glomerular filtration rates (GFR) decrease as a result of hypothermia. In essence, hypothermia increases preglomerular vasoconstriction, thus decreasing both renal blood flow (RBF) and GFR.
Diagnosis
Accurate determination of core temperature often requires a special low temperature thermometer, as most clinical thermometers do not measure accurately below . A low temperature thermometer can be placed in the rectum, esophagus or bladder. Esophageal measurements are the most accurate and are recommended once a person is intubated. Other methods of measurement such as in the mouth, under the arm, or using an infrared ear thermometer are often not accurate.
As a hypothermic person's heart rate may be very slow, prolonged feeling for a pulse could be required before detecting. In 2005, the American Heart Association recommended at least 30–45 seconds to verify the absence of a pulse before initiating CPR. Others recommend a 60-second check.
The classical ECG finding of hypothermia is the Osborn J wave. Also, ventricular fibrillation frequently occurs below and asystole below . The Osborn J may look very similar to those of an acute ST elevation myocardial infarction. Thrombolysis as a reaction to the presence of Osborn J waves is not indicated, as it would only worsen the underlying coagulopathy caused by hypothermia.
Prevention
Staying dry and wearing proper clothing help to prevent hypothermia. Synthetic and wool fabrics are superior to cotton as they provide better insulation when wet and dry. Some synthetic fabrics, such as polypropylene and polyester, are used in clothing designed to wick perspiration away from the body, such as liner socks and moisture-wicking undergarments. Clothing should be loose fitting, as tight clothing reduces the circulation of warm blood. In planning outdoor activity, prepare appropriately for possible cold weather. Those who drink alcohol before or during outdoor activity should ensure at least one sober person responsible for safety is present.
Covering the head is effective, but no more effective than covering any other part of the body. While common folklore says that people lose most of their heat through their heads, heat loss from the head is no more significant than that from other uncovered parts of the body. However, heat loss from the head is significant in infants, whose head is larger relative to the rest of the body than in adults. Several studies have shown that for uncovered infants, lined hats significantly reduce heat loss and thermal stress. Children have a larger surface area per unit mass, and other things being equal should have one more layer of clothing than adults in similar conditions, and the time they spend in cold environments should be limited. However, children are often more active than adults, and may generate more heat. In both adults and children, overexertion causes sweating and thus increases heat loss.
Building a shelter can aid survival where there is danger of death from exposure. Shelters can be constructed out of a variety of materials. Metal can conduct heat away from the occupants and is sometimes best avoided. The shelter should not be too big so body warmth stays near the occupants. Good ventilation is essential especially if a fire will be lit in the shelter. Fires should be put out before the occupants sleep to prevent carbon monoxide poisoning. People caught in very cold, snowy conditions can build an igloo or snow cave to shelter.
The United States Coast Guard promotes using life vests to protect against hypothermia through the 50/50/50 rule: If someone is in water for 50 minutes, they have a 50 percent better chance of survival if they are wearing a life jacket. A heat escape lessening position can be used to increase survival in cold water.
Babies should sleep at and housebound people should be checked regularly to make sure the temperature of the home is at least .
Management
Aggressiveness of treatment is matched to the degree of hypothermia. Treatment ranges from noninvasive, passive external warming to active external rewarming, to active core rewarming. In severe cases resuscitation begins with simultaneous removal from the cold environment and management of the airway, breathing, and circulation. Rapid rewarming is then commenced. Moving the person as little and as gently as possible is recommended as aggressive handling may increase risks of a dysrhythmia.
Hypoglycemia is a frequent complication and needs to be tested for and treated. Intravenous thiamine and glucose is often recommended, as many causes of hypothermia are complicated by Wernicke's encephalopathy.
The UK National Health Service advises against putting a person in a hot bath, massaging their arms and legs, using a heating pad, or giving them alcohol. These measures can cause a rapid fall in blood pressure and potential cardiac arrest.
Rewarming
Rewarming can be done with a number of methods including passive external rewarming, active external rewarming, and active internal rewarming. Passive external rewarming involves the use of a person's own ability to generate heat by providing properly insulated dry clothing and moving to a warm environment. Passive external rewarming is recommended for those with mild hypothermia.
Active external rewarming involves applying warming devices externally, such as a heating blanket. These may function by warmed forced air (Bair Hugger is a commonly used device), chemical reactions, or electricity. In wilderness environments, hypothermia may be helped by placing hot water bottles in both armpits and in the groin. Active external rewarming is recommended for moderate hypothermia. Active core rewarming involves the use of intravenous warmed fluids, irrigation of body cavities with warmed fluids (the chest or abdomen), use of warm humidified inhaled air, or use of extracorporeal rewarming such as via a heart lung machine or extracorporeal membrane oxygenation (ECMO). Extracorporeal rewarming is the fastest method for those with severe hypothermia. When severe hypothermia has led to cardiac arrest, effective extracorporeal warming results in survival with normal mental function about 50% of the time. Chest irrigation is recommended if bypass or ECMO is not possible.
Rewarming shock (or rewarming collapse) is a sudden drop in blood pressure in combination with a low cardiac output which may occur during active treatment of a severely hypothermic person. There was a theoretical concern that external rewarming rather than internal rewarming may increase the risk. These concerns were partly believed to be due to afterdrop, a situation detected during laboratory experiments where there is a continued decrease in core temperature after rewarming has been started. Recent studies have not supported these concerns, and problems are not found with active external rewarming.
Fluids
For people who are alert and able to swallow, drinking warm (not hot) sweetened liquids can help raise the temperature. General medical consensus advises against alcohol and caffeinated drinks. As most hypothermic people are moderately dehydrated due to cold-induced diuresis, warmed intravenous fluids to a temperature of are often recommended.
Cardiac arrest
In those without signs of life, cardiopulmonary resuscitation (CPR) should be continued during active rewarming. For ventricular fibrillation or ventricular tachycardia, a single defibrillation should be attempted. However, people with severe hypothermia may not respond to pacing or defibrillation. It is not known if further defibrillation should be withheld until the core temperature reaches . In Europe, epinephrine is not recommended until the person's core temperature reaches , while the American Heart Association recommends up to three doses of epinephrine before a core temperature of is reached. Once a temperature of has been reached, normal ACLS protocols should be followed.
Prognosis
It is usually recommended not to declare a person dead until their body is warmed to a near normal body temperature of greater than , since extreme hypothermia can suppress heart and brain function. This is summarized in the common saying "You're not dead until you're warm and dead." Exceptions include if there are obvious fatal injuries or the chest is frozen so that it cannot be compressed. If a person was buried in an avalanche for more than 35 minutes and is found with a mouth packed full of snow without a pulse, stopping early may also be reasonable. This is also the case if a person's blood potassium is greater than 12 mmol/L.
Those who are stiff with pupils that do not move may survive if treated aggressively. Survival with good function also occasionally occurs even after the need for hours of CPR. Children who have near-drowning accidents in water near can occasionally be revived, even over an hour after losing consciousness. The cold water lowers the metabolism, allowing the brain to withstand a much longer period of hypoxia. While survival is possible, mortality from severe or profound hypothermia remains high despite optimal treatment. Studies estimate mortality at between 38% and 75%.
In those who have hypothermia due to another underlying health problem, when death occurs it is frequently from that underlying health problem.
Epidemiology
Between 1995 and 2004 in the United States, an average of 1,560 cold-related emergency department visits occurred per year and in the years 1999 to 2004, an average of 647 people died per year due to hypothermia. Of deaths reported between 1999 and 2002 in the US, 49% of those affected were 65 years or older and two-thirds were male. Most deaths were not work related (63%) and 23% of affected people were at home. Hypothermia was most common during the autumn and winter months of October through March. In the United Kingdom, an estimated 300 deaths per year are due to hypothermia, whereas the annual incidence of hypothermia-related deaths in Canada is 8,000.
History
Hypothermia has played a major role in the success or failure of many military campaigns, from Hannibal's loss of nearly half his men in the Second Punic War (218 B.C.) to the near destruction of Napoleon's armies in Russia in 1812. Men wandered around confused by hypothermia, some lost consciousness and died, others shivered, later developed torpor, and tended to sleep. Others too weak to walk fell on their knees; some stayed that way for some time resisting death. The pulse of some was weak and hard to detect; others groaned; yet others had eyes open and wild with quiet delirium. Deaths from hypothermia in Russian regions continued through the first and second world wars, especially on the Eastern Front in battles such as the Battle of Moscow and the Battle of Stalingrad where German soldiers were not provided with winter clothing.
Civilian examples of deaths caused by hypothermia occurred during the sinkings of the RMS Titanic and RMS Lusitania, and more recently of the MS Estonia.
Antarctic explorers developed hypothermia; Ernest Shackleton and his team measured body temperatures "below 94.2°, which spells death at home", though this probably referred to oral temperatures rather than core temperature and corresponded to mild hypothermia. One of Scott's team, Atkinson, became confused through hypothermia.
Nazi human experimentation during World War II amounting to medical torture included hypothermia experiments, which killed many victims. There were 360 to 400 experiments and 280 to 300 subjects, indicating some had more than one experiment performed on them. Various methods of rewarming were attempted: "One assistant later testified that some victims were thrown into boiling water for rewarming".
In 2024, at least six babies in Gaza died of hypothermia under the harsh rain and cold, which they had to endure in flimsy camps due to the bombing of their homes and forced displacement inflicted upon them by the IDF. Three of the babies died in the coastal zone of Al-Mawasi, which there families had been forced to evacuate to, given its designation as "safe zone".
Medical use
Various degrees of hypothermia may be deliberately induced in medicine for purposes of treatment of brain injury, or lowering metabolism so that total brain ischemia can be tolerated for a short time. Deep hypothermic circulatory arrest is a medical technique in which the brain is cooled as low as 10 °C, which allows the heart to be stopped and blood pressure to be lowered to zero, for the treatment of aneurysms and other circulatory problems that do not tolerate arterial pressure or blood flow. The time limit for this technique, as also for accidental arrest in ice water (which internal temperatures may drop to as low as 15 °C), is about one hour.
Other animals
Hypothermia can happen in most mammals in cold weather and can be fatal. Baby mammals such as kittens are unable to regulate their body temperatures and have a risk of hypothermia if they are not kept warm by their mothers.
Many animals other than humans often induce hypothermia during hibernation or torpor.
Water bears (Tardigrade), microscopic multicellular organisms, can survive freezing at low temperatures by replacing most of their internal water with the sugar trehalose, preventing the crystallization that otherwise damages cell membranes.
| Biology and health sciences | Injury | null |
146924 | https://en.wikipedia.org/wiki/Radial%20engine | Radial engine | The radial engine is a reciprocating type internal combustion engine configuration in which the cylinders "radiate" outward from a central crankcase like the spokes of a wheel. It resembles a stylized star when viewed from the front, and is called a "star engine" in some other languages.
The radial configuration was commonly used for aircraft engines before gas turbine engines became predominant.
Engine operation
Since the axes of the cylinders are coplanar, the connecting rods cannot all be directly attached to the crankshaft unless mechanically complex forked connecting rods are used, none of which have been successful. Instead, the pistons are connected to the crankshaft with a master-and-articulating-rod assembly. One piston, the uppermost one in the animation, has a master rod with a direct attachment to the crankshaft. The remaining pistons pin their connecting rods' attachments to rings around the edge of the master rod. Extra "rows" of radial cylinders can be added in order to increase the capacity of the engine without adding to its diameter.
Four-stroke radials have an odd number of cylinders per row, so that a consistent every-other-piston firing order can be maintained, providing smooth operation. For example, on a five-cylinder engine the firing order is 1, 3, 5, 2, 4, and back to cylinder 1. Moreover, this always leaves a one-piston gap between the piston on its combustion stroke and the piston on compression. The active stroke directly helps compress the next cylinder to fire, making the motion more uniform. If an even number of cylinders were used, an equally timed firing cycle would not be feasible.
As with most four-strokes, the crankshaft takes two revolutions to complete the four strokes of each piston (intake, compression, combustion, exhaust). The camshaft ring is geared to spin slower and in the opposite direction to the crankshaft. Its cam lobes are placed in two rows; one for the intake valves and one for the exhaust valves. The radial engine normally uses fewer cam lobes than other types. For example, in the engine in the animated illustration, four cam lobes serve all 10 valves across the five cylinders, whereas 10 would be required for a typical inline engine with the same number of cylinders and valves.
Most radial engines use overhead poppet valves driven by pushrods and lifters on a cam plate which is concentric with the crankshaft, with a few smaller radials, like the Kinner B-5 and Russian Shvetsov M-11, using individual camshafts within the crankcase for each cylinder. A few engines use sleeve valves such as the 14-cylinder Bristol Hercules and the 18-cylinder Bristol Centaurus, which are quieter and smoother running but require much tighter manufacturing tolerances.
History
Aircraft
C. M. Manly constructed a water-cooled five-cylinder radial engine in 1901, a conversion of one of Stephen Balzer's rotary engines, for Langley's Aerodrome aircraft. Manly's engine produced at 950 rpm.
In 1903–1904 Jacob Ellehammer used his experience constructing motorcycles to build the world's first air-cooled radial engine, a three-cylinder engine which he used as the basis for a more powerful five-cylinder model in 1907. This was installed in his triplane and made a number of short free-flight hops.
Another early radial engine was the three-cylinder Anzani, originally built as a W3 "fan" configuration, one of which powered Louis Blériot's Blériot XI across the English Channel. Before 1914, Alessandro Anzani had developed radial engines ranging from 3 cylinders (spaced 120° apart) — early enough to have been used on a few French-built examples of the famous Blériot XI from the original Blériot factory — to a massive 20-cylinder engine of , with its cylinders arranged in four rows of five cylinders apiece.
Most radial engines are air-cooled, but one of the most successful of the early radial engines (and the earliest "stationary" design produced for World War I combat aircraft) was the Salmson 9Z series of nine-cylinder water-cooled radial engines that were produced in large numbers. Georges Canton and Pierre Unné patented the original engine design in 1909, offering it to the Salmson company; the engine was often known as the Canton-Unné.
From 1909 to 1919 the radial engine was overshadowed by its close relative, the rotary engine, which differed from the so-called "stationary" radial in that the crankcase and cylinders revolved with the propeller. It was similar in concept to the later radial, the main difference being that the propeller was bolted to the engine, and the crankshaft to the airframe. The problem of the cooling of the cylinders, a major factor with the early "stationary" radials, was alleviated by the engine generating its own cooling airflow.
In World War I many French and other Allied aircraft flew with Gnome, Le Rhône, Clerget, and Bentley rotary engines, the ultimate examples of which reached although none of those over were successful. By 1917 rotary engine development was lagging behind new inline and V-type engines, which by 1918 were producing as much as , and were powering almost all of the new French and British combat aircraft.
Most German aircraft of the time used water-cooled inline 6-cylinder engines. Motorenfabrik Oberursel made licensed copies of the Gnome and Le Rhône rotary powerplants, and Siemens-Halske built their own designs, including the Siemens-Halske Sh.III eleven-cylinder rotary engine, which was unusual for the period in being geared through a bevel geartrain in the rear end of the crankcase without the crankshaft being firmly mounted to the aircraft's airframe, so that the engine's internal working components (fully internal crankshaft "floating" in its crankcase bearings, with its conrods and pistons) were spun in the opposing direction to the crankcase and cylinders, which still rotated as the propeller itself did since it was still firmly fastened to the crankcase's frontside, as with regular umlaufmotor German rotaries.
By the end of the war the rotary engine had reached the limits of the design, particularly in regard to the amount of fuel and air that could be drawn into the cylinders through the hollow crankshaft, while advances in both metallurgy and cylinder cooling finally allowed stationary radial engines to supersede rotary engines. In the early 1920s Le Rhône converted a number of their rotary engines into stationary radial engines.
By 1918 the potential advantages of air-cooled radials over the water-cooled inline engine and air-cooled rotary engine that had powered World War I aircraft were appreciated but were unrealized. British designers had produced the ABC Dragonfly radial in 1917, but were unable to resolve the cooling problems, and it was not until the 1920s that Bristol and Armstrong Siddeley produced reliable air-cooled radials such as the Bristol Jupiter and the Armstrong Siddeley Jaguar.
In the United States the National Advisory Committee for Aeronautics (NACA) noted in 1920 that air-cooled radials could offer an increase in power-to-weight ratio and reliability; by 1921 the U.S. Navy had announced it would only order aircraft fitted with air-cooled radials and other naval air arms followed suit. Charles Lawrance's J-1 engine was developed in 1922 with Navy funding, and using aluminum cylinders with steel liners ran for an unprecedented 300 hours, at a time when 50 hours endurance was normal. At the urging of the Army and Navy the Wright Aeronautical Corporation bought Lawrance's company, and subsequent engines were built under the Wright name. The radial engines gave confidence to Navy pilots performing long-range overwater flights.
Wright's J-5 Whirlwind radial engine of 1925 was widely claimed as "the first truly reliable aircraft engine". Wright employed Giuseppe Mario Bellanca to design an aircraft to showcase it, and the result was the Wright-Bellanca WB-1, which first flew later that year. The J-5 was used on many advanced aircraft of the day, including Charles Lindbergh's Spirit of St. Louis, in which he made the first solo trans-Atlantic flight.
In 1925 the American Pratt & Whitney company was founded, competing with Wright's radial engines. Pratt & Whitney's initial offering, the R-1340 Wasp, was test run later that year, beginning a line of engines over the next 25 years that included the 14-cylinder, twin-row Pratt & Whitney R-1830 Twin Wasp. More Twin Wasps were produced than any other aviation piston engine in the history of aviation; nearly 175,000 were built.
In the United Kingdom the Bristol Aeroplane Company was concentrating on developing radials such as the Jupiter, Mercury, and sleeve valve Hercules radials. Germany, Japan, and the Soviet Union started with building licensed versions of the Armstrong Siddeley, Bristol, Wright, or Pratt & Whitney radials before producing their own improved versions. France continued its development of various rotary engines but also produced engines derived from Bristol designs, especially the Jupiter.
Although other piston configurations and turboprops have taken over in modern propeller-driven aircraft, Rare Bear, which is a Grumman F8F Bearcat equipped with a Wright R-3350 Duplex-Cyclone radial engine, is still the fastest piston-powered aircraft.
125,334 of the American twin-row, 18-cylinder Pratt & Whitney R-2800 Double Wasp, with a displacement of 2,800 in3 (46 L) and between 2,000 and 2,400 hp (1,500-1,800 kW), powered the American single-engine Vought F4U Corsair, Grumman F6F Hellcat, Republic P-47 Thunderbolt, twin-engine Martin B-26 Marauder, Douglas A-26 Invader, Northrop P-61 Black Widow, etc. The same firm's aforementioned smaller-displacement (at 30 litres), Twin Wasp 14-cylinder twin-row radial was used as the main engine design for the B-24 Liberator, PBY Catalina, and Douglas C-47, each design being among the production leaders in all-time production numbers for each type of airframe design.
The American Wright Cyclone series twin-row radials powered American warplanes: the nearly-43 litre displacement, 14-cylinder Twin Cyclone powered the single-engine Grumman TBF Avenger, twin-engine North American B-25 Mitchell, and some versions of the Douglas A-20 Havoc, with the massive twin-row, nearly 55-litre displacement, 18-cylinder Duplex-Cyclone powering the four-engine Boeing B-29 Superfortress and others.
The Soviet Shvetsov OKB-19 design bureau was the sole source of design for all of the Soviet government factory-produced radial engines used in its World War II aircraft, starting with the Shvetsov M-25 (itself based on the American Wright Cyclone 9's design) and going on to design the 41-litre displacement Shvetsov ASh-82 fourteen cylinder radial for fighters, and the massive, 58-litre displacement Shvetsov ASh-73 eighteen-cylinder radial in 1946 - the smallest-displacement radial design from the Shvetsov OKB during the war was the indigenously designed, 8.6 litre displacement Shvetsov M-11 five cylinder radial.
Over 28,000 of the German 42-litre displacement, 14-cylinder, two-row BMW 801, with between 1,560 and 2,000 PS (1,540-1,970 hp, or 1,150-1,470 kW), powered the German single-seat, single-engine Focke-Wulf Fw 190 Würger, and twin-engine Junkers Ju 88.
In Japan, most airplanes were powered by air-cooled radial engines like the 14-cylinder Mitsubishi Zuisei (11,903 units, e.g. Kawasaki Ki-45), Mitsubishi Kinsei (12,228 units, e.g. Aichi D3A), Mitsubishi Kasei (16,486 units, e.g. Kawanishi H8K), Nakajima Sakae (30,233 units, e.g. Mitsubishi A6M and Nakajima Ki-43), and 18-cylinder Nakajima Homare (9,089 units, e.g. Nakajima Ki-84). The Kawasaki Ki-61 and Yokosuka D4Y were rare examples of Japanese liquid-cooled inline engine aircraft at that time but later, they were also redesigned to fit radial engines as the Kawasaki Ki-100 and Yokosuka D4Y3.
In Britain, Bristol produced both sleeve valved and conventional poppet valved radials: of the sleeve valved designs, more than 57,400 Hercules engines powered the Vickers Wellington, Short Stirling, Handley Page Halifax, and some versions of the Avro Lancaster, over 8,000 of the pioneering sleeve-valved Bristol Perseus were used in various types, and more than 2,500 of the largest-displacement production British radial from the Bristol firm to use sleeve valving, the Bristol Centaurus were used to power the Hawker Tempest II and Sea Fury. The same firm's poppet-valved radials included: around 32,000 of Bristol Pegasus used in the Short Sunderland, Handley Page Hampden, and Fairey Swordfish and over 20,000 examples of the firm's 1925-origin nine-cylinder Mercury were used to power the Westland Lysander, Bristol Blenheim, and Blackburn Skua.
Tanks
In the years leading up to World War II, as the need for armored vehicles was realized, designers were faced with the problem of how to power the vehicles, and turned to using aircraft engines, among them radial types. The radial aircraft engines provided greater power-to-weight ratios and were more reliable than conventional inline vehicle engines available at the time. This reliance had a downside though: if the engines were mounted vertically, as in the M3 Lee and M4 Sherman, their comparatively large diameter gave the tank a higher silhouette than designs using inline engines.
The Continental R-670, a 7-cylinder radial aero engine which first flew in 1931, became a widely used tank powerplant, being installed in the M1 Combat Car, M2 Light Tank, M3 Stuart, M3 Lee, and LVT-2 Water Buffalo.
The Guiberson T-1020, a 9-cylinder radial diesel aero engine, was used in the M1A1E1, while the Continental R975 saw service in the M4 Sherman, M7 Priest, M18 Hellcat tank destroyer, and the M44 self propelled howitzer.
Modern radials
A number of companies continue to build radials today. Vedeneyev produces the M-14P radial of as used on Yakovlev and Sukhoi aerobatic aircraft. The M-14P is also used by builders of homebuilt aircraft, such as the Culp Special, and Culp Sopwith Pup, Pitts S12 "Monster" and the Murphy "Moose". 7-cylinder and 9-cylinder engines are available from Australia's Rotec Aerosport. HCI Aviation offers the R180 5-cylinder () and R220 7-cylinder (), available "ready to fly" and as a build-it-yourself kit. Verner Motor of the Czech Republic builds several radial engines ranging in power from . Miniature radial engines for model airplanes are available from O. S. Engines, Saito Seisakusho of Japan, and Shijiazhuang of China, and Evolution (designed by Wolfgang Seidel of Germany, and made in India) and Technopower in the US.
Comparison with inline engines
Liquid cooling systems are generally more vulnerable to battle damage. Even minor shrapnel damage can easily result in a loss of coolant and consequent engine overheating, while an air-cooled radial engine may be largely unaffected by minor damage. Radials have shorter and stiffer crankshafts, a single-bank radial engine needing only two crankshaft bearings as opposed to the seven required for a liquid-cooled, six-cylinder, inline engine of similar stiffness.
While a single-bank radial permits all cylinders to be cooled equally, the same is not true for multi-row engines where the rear cylinders can be affected by the heat coming off the front row, and air flow being masked.
A potential disadvantage of radial engines is that having the cylinders exposed to the airflow increases drag considerably. The answer was the addition of specially designed cowlings with baffles to force the air between the cylinders. The first effective drag-reducing cowling that didn't impair engine cooling was the British Townend ring or "drag ring" which formed a narrow band around the engine covering the cylinder heads, reducing drag. The National Advisory Committee for Aeronautics studied the problem, developing the NACA cowling which further reduced drag and improved cooling. Nearly all aircraft radial engines since have used NACA-type cowlings.{{refn|group=Note|It has been claimed that the NACA cowling generated extra thrust due to the Meredith Effect, whereby the heat added to the air being forced through the ducts between the cylinders expanded the exhausting cooling air, producing thrust when forced through a nozzle. The Meredith effect requires high airspeed and careful design to generate a suitable high speed exhaust of the heated air – the NACA cowling was not designed to achieve this, nor would the effect have been significant at low airspeeds.<ref name=becker>Becker, J.; [http://www.hq.nasa.gov/pao/History/SP-445/ch5-5.htm The high-speed frontier: Case histories of four NACA programs, 1920- SP-445, NASA (1980), Chapter 5: High-speed Cowlings, Air Inlets and Outlets, and Internal-Flow Systems: The ramjet investigation]</ref> The effect was put to use in the radiators of several mid-1940s aircraft that used liquid-cooled engines such as the Spitfire and Mustang, and it offered a minor improvement in later radial-engined aircraft, including the Fw 190.}}
While inline liquid-cooled engines continued to be common in new designs until late in World War II, radial engines dominated afterwards until overtaken by jet engines, with the late-war Hawker Sea Fury and Grumman F8F Bearcat, two of the fastest production piston-engined aircraft ever built, using radial engines.
Hydrolock
Whenever a radial engine remains shut down for more than a few minutes, oil or fuel may drain into the combustion chambers of the lower cylinders or accumulate in the lower intake pipes, ready to be drawn into the cylinders when the engine starts. As the piston approaches top dead center (TDC) of the compression stroke, this liquid, being incompressible, stops piston movement. Starting or attempting to start the engine in such condition may result in a bent or broken connecting rod.
Other types of radial engine
Multi-row radials
Originally radial engines had one row of cylinders, but as engine sizes increased it became necessary to add extra rows. The first radial-configuration engine known to use a twin-row design was the 160 hp Gnôme "Double Lambda" rotary engine of 1912, designed as a 14-cylinder twin-row version of the firm's 80 hp Lambda single-row seven-cylinder rotary, however reliability and cooling problems limited its success.
Two-row designs began to appear in large numbers during the 1930s, when aircraft size and weight grew to the point where single-row engines of the required power were simply too large to be practical. Two-row designs often had cooling problems with the rear bank of cylinders, but a variety of baffles and fins were introduced that largely eliminated these problems. The downside was a relatively large frontal area that had to be left open to provide enough airflow, which increased drag. This led to significant arguments in the industry in the late 1930s about the possibility of using radials for high-speed aircraft like modern fighters.
The solution was introduced with the BMW 801 14-cylinder twin-row radial. Kurt Tank designed a new cooling system for this engine that used a high-speed fan to blow compressed air into channels that carry air to the middle of the banks, where a series of baffles directed the air over all of the cylinders. This allowed the cowling to be tightly fitted around the engine, reducing drag, while still providing (after a number of experiments and modifications) enough cooling air to the rear. This basic concept was soon copied by many other manufacturers, and many late-WWII aircraft returned to the radial design as newer and much larger designs began to be introduced. Examples include the Bristol Centaurus in the Hawker Sea Fury, and the Shvetsov ASh-82 in the Lavochkin La-7.
For even greater power, adding further rows was not considered viable due to the difficulty of providing the required airflow to the rear banks. Larger engines were designed, mostly using water cooling although this greatly increased complexity and eliminated some of the advantages of the radial air-cooled design. One example of this concept is the BMW 803, which never entered service.
A major study into the airflow around radials using wind tunnels and other systems was carried out in the US, and demonstrated that ample airflow was available with careful design. This led to the R-4360, which has 28 cylinders arranged in a 4 row corncob'' configuration. The R-4360 saw service on large American aircraft in the post-World War II period. The US and Soviet Union continued experiments with larger radials, but the UK abandoned such designs in favour of newer versions of the Centaurus and rapid movement to the use of turboprops such as the Armstrong Siddeley Python and Bristol Proteus, which easily produced more power than radials without the weight or complexity.
Large radials continued to be built for other uses, although they are no longer common. An example is the 5-ton Zvezda M503 diesel engine with 42 cylinders in 6 rows of 7, displacing and producing . Three of these were used on the fast Osa class missile boats. Another one was the Lycoming XR-7755 which was the largest piston aircraft engine ever built in the United States with 36 cylinders totaling about 7,750 in3 (127 L) of displacement and a power output of 5,000 horsepower (3,700 kilowatts).
Diesel radials
While most radial engines have been produced for gasoline, there have been diesel radial engines. Two major advantages favour diesel engines — lower fuel consumption and reduced fire risk.
Packard
Packard designed and built a 9-cylinder 980 cubic inch (16.06 litre) displacement diesel radial aircraft engine, the DR-980, in 1928. On 28 May 1931, a DR-980 powered Bellanca CH-300, with 481 gallons of fuel, piloted by Walter Edwin Lees and Frederick Brossy set a record for staying aloft for 84 hours and 32 minutes without being refueled. This record stood for 55 years until broken by the Rutan Voyager.
Bristol
The experimental Bristol Phoenix of 1928–1932 was successfully flight tested in a Westland Wapiti and set altitude records in 1934 that lasted until World War II.
Clerget
In 1932 the French company Clerget developed the 14D, a 14-cylinder two-stroke diesel radial engine. After a series of improvements, in 1938 the 14F2 model produced at 1910 rpm cruise power, with a power-to-weight ratio near that of contemporary gasoline engines and a specific fuel consumption of roughly 80% that for an equivalent gasoline engine. During WWII the research continued, but no mass-production occurred because of the Nazi occupation. By 1943 the engine had grown to produce over with a turbocharger. After the war, the Clerget company was integrated in the SNECMA company and had plans for a 32-cylinder diesel engine of , but in 1947 the company abandoned piston engine development in favour of the emerging turbine engines.
Nordberg
The Nordberg Manufacturing Company of the United States developed and produced a series of large two-stroke radial diesel engines from the late 1940s for electrical production, primarily at aluminum smelters and for pumping water. They differed from most radials in that they had an even number of cylinders in a single bank (or row) and an unusual double master connecting rod. Variants were built that could be run on either diesel oil or gasoline or mixtures of both. A number of powerhouse installations utilising large numbers of these engines were made in the U.S.
EMD
Electro-Motive Diesel (EMD) built the "pancake" engines 16-184 and 16-338 for marine use.
Zoche
Zoche aero-diesels are a prototype radial design that have an even number of cylinders, either four or eight; but this is not problematic, because they are two-stroke engines, with twice the number of power strokes as a four-stroke engine per crankshaft rotation.
Compressed air radial engines
A number of radial motors operating on compressed air have been designed, mostly for use in model airplanes and in gas compressors.
Model radial engines
A number of multi-cylinder 4-stroke model engines have been commercially available in a radial configuration, beginning with the Japanese O.S. Max firm's FR5-300 five-cylinder, 3.0 cu.in. (50 cm3) displacement "Sirius" radial in 1986. The American "Technopower" firm had made smaller-displacement five- and seven-cylinder model radial engines as early as 1976, but the OS firm's engine was the first mass-produced radial engine design in aeromodelling history. The rival Saito Seisakusho firm in Japan has since produced a similarly sized five-cylinder radial four-stroke model engine of their own as a direct rival to the OS design, with Saito also creating a series of three-cylinder methanol and gasoline-fueled model radial engines ranging from 0.90 cu.in. (15 cm3) to 4.50 cu.in. (75 cm3) in displacement, also all now available in spark-ignition format up to 84 cm3 displacement for use with gasoline. The German Seidel firm formerly made both seven- and nine-cylinder "large" (starting at 35 cm3 displacement) radio control model radial engines, mostly for glow plug ignition, with an experimental fourteen-cylinder twin-row radial being tried out - the American Evolution firm now sells the Seidel-designed radials, with their manufacturing being done in India.
| Technology | Aircraft components | null |
146977 | https://en.wikipedia.org/wiki/DVB | DVB | Digital Video Broadcasting (DVB) is a set of international open standards for digital television. DVB standards are maintained by the DVB Project, an international industry consortium, and are published by a Joint Technical Committee (JTC) of the European Telecommunications Standards Institute (ETSI), European Committee for Electrotechnical Standardization (CENELEC) and European Broadcasting Union (EBU).
Transmission
DVB systems distribute data using a variety of approaches, including:
Satellite: DVB-S, DVB-DSNG, DVB-S2, DVB-S2X and DVB-SH
DVB-SMATV for distribution via SMATV
Cable: DVB-C, DVB-C2
Terrestrial television: DVB-T, DVB-T2
Digital terrestrial television for handhelds: DVB-H, DVB-SH
Microwave: using DTT (DVB-MT), the MMDS (DVB-MC), and/or MVDS standards (DVB-MS)
These standards define the physical layer and data link layer of the distribution system. Devices interact with the physical layer via a synchronous parallel interface (SPI), synchronous serial interface (SSI) or asynchronous serial interface (ASI). All data is transmitted in MPEG transport streams with some additional constraints (DVB-MPEG). A standard for temporally-compressed distribution to mobile devices (DVB-H) was published in November 2004.
These distribution systems differ mainly in the modulation schemes used and error correcting codes used, due to the different technical constraints. DVB-S (SHF) uses QPSK, 8-PSK or 16-QAM. DVB-S2 uses QPSK, 8-PSK, 16-APSK or 32-APSK, at the broadcasters decision. QPSK and 8-PSK are the only versions regularly used. DVB-C (VHF/UHF) uses QAM: 16-QAM, 32-QAM, 64-QAM, 128-QAM or 256-QAM. Lastly, DVB-T (VHF/UHF) uses 16-QAM or 64-QAM (or QPSK) in combination with (C)OFDM and can support hierarchical modulation.
The DVB-T2 specification was approved by the DVB Steering Board in June 2008 and sent to ETSI for adoption as a formal standard. ETSI adopted the standard on 9 September 2009. The DVB-T2 standard gives more robust TV reception and increases the possible bit rate by over 30% for single transmitters (as in the UK) and will increase the maximum bit rate by over 50% in large single-frequency networks (as in Germany and Sweden).
DVB has established a 3D TV group (CM-3DTV) to identify "what kind of 3D-TV solution does the market want and need, and how can DVB play an active part in the creation of that solution?" The CM-3DTV group held a DVB 3D-TV Kick-off Workshop in Geneva on 25 January 2010, followed by the first CM-3DTV meeting the next day. DVB now defines a new standard for 3D video broadcast: DVB 3D-TV.
Modes and features of latest DVB-x2 system standards in comparison:
Content
Digital video content is encoded using discrete cosine transform (DCT) based video coding standards, such as the H.26x and MPEG formats. Digital audio content is encoded using modified discrete cosine transform (MDCT) based audio coding standards, such as Advanced Audio Coding (AAC), Dolby Digital (AC-3) and MP3.
Besides digital audio and digital video transmission, DVB also defines data connections (DVB-DATA - EN 301 192) with return channels (DVB-RC) for several media (DECT, GSM, PSTN/ISDN, satellite etc.) and protocols (DVB-IPTV: Internet Protocol; DVB-NPI: network protocol independent).
Older technologies such as teletext (DVB-TXT) and vertical blanking interval data (DVB-VBI) are also supported by the standards to ease conversion. However, for many applications more advanced alternatives like DVB-SUB for subtitling are available.
Encryption and metadata
The conditional access system (DVB-CA) defines a Common Scrambling Algorithm (DVB-CSA) and a physical Common Interface (DVB-CI) for accessing scrambled content. DVB-CA providers develop their wholly proprietary conditional access systems with reference to these specifications. Multiple simultaneous CA systems can be assigned to a scrambled DVB program stream providing operational and commercial flexibility for the service provider.
The DVB Project developed a Content Protection and Copy Management system for protecting content after it has been received (DVB-CPCM), which was intended to allow flexible use of recorded content on a home network or beyond, while preventing unconstrained sharing on the Internet. DVB-CPCM was the source of much controversy in the popular press and it was said that CPCM was the DVB Project's answer to the failed American Broadcast Flag. The DVB-CPCM specifications, which were standardized by ETSI as a multipart document (TS 102 825) between 2008 and 2013, were deprecated by the DVB Steering Board in February 2019.
DVB transports include metadata called Service Information (DVB-SI, ETSI EN 300 468, ETSI TR 101 211) that links the various elementary streams into coherent programs and provides human-readable descriptions for electronic program guides as well as for automatic searching and filtering. The dating system used with this metadata suffers from a year 2038 problem in which due to the limited 16 bits and modified Julian day offset used will cause an overflow issue similar to the year 2000 problem. By comparison, the rival DigiCipher 2 based ATSC system will not have this issue until 2048 due in part to 32 bits being used.
DVB adopted a profile of the metadata defined by the TV-Anytime Forum (DVB-TVA, ETSI TS 102323). This is an XML Schema based technology and the DVB profile is tailored for enhanced Personal Digital Recorders.
In the early 2000s, DVB started an activity to develop specifications for IPTV (DVB-IPI, ETSI TR 102 033, ETSI TS 102 034, ETSI TS 102 814), which also included metadata definitions for a broadband content guide (DVB-BCG, ETSI TS 102 539).
DVB-I
In October 2017, the DVB Project established a working group to begin the definition of a specification for "standalone TV services over IP, referred to as DVB-I services". Work on the commercial requirements for DVB-I began in January 2018 and the terms of reference were agreed in March 2018.
The DVB-I specification, titled "Service Discovery and Programme Metadata for DVB-I", was approved by the DVB Project in November 2019
and first published as DVB BlueBook A177 in June 2020 and as an ETSI standard TS 103 770 in November 2020.
The DVB-I specification defines ways in which devices and displays connected to the internet can discover and access sets of audiovisual media services. These can include services delivered online through fixed and wireless Internet Protocol connections as well as broadcast radio and television channels received over radio frequency networks using traditional cable, satellite, or terrestrial transmissions.
Tests and pilots of DVB-I services have been undertaken in several countries including Iran, Germany, Italy, Spain and Ireland.
Software platform
The DVB Multimedia Home Platform (DVB-MHP) defines a Java-based platform for the development of consumer video system applications. In addition to providing abstractions for many DVB and MPEG-2 concepts, it provides interfaces for other features like network card control, application download, and layered graphics.
Return channel
DVB has standardized a number of return channels that work together with DVB(-S/T/C) to create bi-directional communication. RCS is short for Return Channel Satellite, and specifies return channels in C, Ku and Ka frequency bands with return bandwidth of up to 2 Mbit/s. DVB-RCT is short for Return Channel Terrestrial, specified by ETSI EN 301958.
Service discovery
The DVB-I standard (ETSI TS 103 770) defines an internet-based request and response mechanism to discover and access audiovisual services delivered over traditional digital broadcast transmissions or Internet Protocol networks and present them in a unified way.
Adoption
DVB-S and DVB-C were ratified in 1994. DVB-T was ratified in early 1997. The first commercial DVB-T broadcasts were performed by the United Kingdom's Digital TV Group in late 1998. In 2003 Berlin, Germany was the first area to completely stop broadcasting analogue TV signals. Most European countries are fully covered by digital television and many have switched off PAL/SECAM services.
DVB standards are used throughout Europe, as well as in Australia, South Africa and India. They are also used for cable and satellite broadcasting in most Asian, African and many South American countries. Some have chosen ISDB-T instead of DVB-T and a few (United States, Canada, Mexico and South Korea) have chosen ATSC instead of DVB-T.
Africa
Kenya
DVB-T broadcasts were launched by the President of Kenya, Mwai Kibaki on 9 December 2009. Broadcasts are using H.264, with the University of Nairobi supplying the decoders. Kenya has also been broadcasting DVB-H since July 2009, available on selected Nokia and ZTE handsets on the Safaricom and other GSM networks.
Madagascar
Since 2011, the pay TV operator Blueline launched a DVB-T service branded BluelineTV. It supplies both smart cards and set-top-boxes.
South Africa
Since 1995, the pay TV operator DStv used the DVB-S standard to broadcast its services. In 2010, it started a DVB over IP service, and in 2011 it started DStv mobile using the DVB-H standard.
In late 2010, the South African cabinet endorsed a decision by a Southern African Development Community (SADC) task team to adopt the DVB-T2 standard.
Asia
Hong Kong
In Hong Kong, several cable TV operators such as TVB Pay Vision and Cable TV have already started using DVB-S or DVB-C. The government however has adopted the DMB-T/H standard, developed in mainland China, for its digital terrestrial broadcasting services which has started since 31 December 2007.
Iran
On 17 March 2009, DVB-H and DVB-T H.264/AAC broadcasting started in Tehran by the IRIB. DVB-T broadcasting is now widely available in other cities such as Isfahan, Mashhad, Shiraz, Qom, Tabriz and Rasht as well.
Israel
DVB-T broadcasts using H.264 commenced in Israel on 1 June 2009 with the broadcast trial and the full broadcast began on 2 August 2009. Analog broadcasts were originally planned to end in 18 months after the launch, but analog broadcasts were switched off on 31 March 2011 instead.
During 2010, DVB-T broadcasts have become widely available in most of Israel and an EPG was added to the broadcasts.
Japan
With the exception of SKY PerfecTV!, Japan uses different formats in all areas (ISDB), which are however quite similar to their DVB counterparts. SkyPerfect is a satellite provider using DVB on its 124 and 128 degrees east satellites. Its satellite at 110 degrees east does not use DVB, however.
Malaysia
In Malaysia, a new pay television station MiTV began service in September 2005 using DVB-IPTV technology while lone satellite programming provider ASTRO has been transmitting in DVB-S since its inception in 1996. Free-to-air DVB-T trials began in late 2006 with a simulcast of both TV1 and TV2 plus a new channel called RTM3/RTMi. In April 2007, RTM announced that the outcome of the test was favourable and that it expected DVB-T to go public by the end of 2007. However, the system did not go public as planned. As of 2008, the trial digital line-up has expanded to include a music television channel called Muzik Aktif, and a sports channel called Arena, with a news channel called Berita Aktif planned for inclusion in the extended trials soon. Also, high definition trials were performed during the Beijing Olympics and the outcome was also favourable. It was announced that the system would go public in 2009.
In 2009, MiTV closed down, changed its name to U-Television and announced that it was changing to scrambled DVB-T upon relaunch instead of the DVB-IPTV system used prior to shutting down. However, RTM's digital network again did not go public, although around this time TVs that are first-generation DVB-T capable went on sale. The government has since announced that they will be deploying DVB-T2 instead in stages starting in mid-2015 and analog shutoff has been delayed to April 2019.
Philippines
In the Philippines, DVB-S and DVB-S2 are the two broadcast standards currently used by satellite companies, while DVB-C is also used by some cable companies. The government adopted DVB-T in November 2006 for digital terrestrial broadcasting but a year later, it considered other standards to replace DVB-T. The country has chosen the ISDB-T system instead of DVB-T.
Taiwan
In Taiwan, some digital cable television systems use DVB-C, though most customers still use analogue NTSC cable television. The government planned adopting ATSC or the Japanese ISDB-T standard as NTSC's replacement. However, the country has chosen the European DVB-T system instead. Public Television Service (PTS) and Formosan TV provide high definition television. The former has the channel HiHD; the latter uses its HD channel for broadcasting MLB baseball.
Europe
Cyprus
Cyprus uses DVB-T with MPEG-4 encoding. Analogue transmission stopped on 1 July 2011 for all channels except CyBC 1.
Denmark
In Denmark, DVB-T replaced the analog transmission system for TV on 1 November 2009. Danish national digital TV transmission has been outsourced to the company Boxer TV A/S, acting as gatekeeper organization for terrestrial TV transmission in Denmark. However, there are still several free channels from DR.
Finland
DVB-T transmissions were launched on 21 August 2001. The analogue networks continued alongside the digital ones until 1 September 2007, when they were shut down nationwide. Before the analogue switchoff, the terrestrial network had three multiplexes: MUX A, MUX B and MUX C. MUX A contained the channels of the public broadcaster Yleisradio and MUX B was shared between the two commercial broadcasters: MTV3 and Nelonen. MUX C contained channels of various other broadcasters. After the analogue closedown, a fourth multiplex named MUX E was launched. All of the Yleisradio (YLE) channels are broadcast free-to-air, likewise a handful of commercial ones including MTV3, Nelonen, Subtv, Jim, Nelonen Sport, Liv, FOX, TV5 Finland, AVA and Kutonen. There are also several pay channels sold by PlusTV.
Italy
In Italy, DVB-S started in 1996 and the final analogue broadcasts were terminated in 2005. The switch-off from analogue terrestrial network to DVB-T started on 15 October 2008. Analogue broadcast was ended on 4 July 2012 after nearly four years of transition in phases.
Netherlands
In the Netherlands, DVB-S broadcasting started on 1 July 1996, satellite provider MultiChoice (now CanalDigitaal) switched off the analogue service shortly after on 18 August 1996. DVB-T broadcasting started April 2003, and terrestrial analog broadcasting was switched off December 2006. It was initially marketed by Digitenne but later by KPN. Multiplex 1 contains the NPO 1, NPO 2 and NPO 3 national TV channels, and a regional channel. Multiplexes 2~5 have the other encrypted commercial and international channels. Multiplex 1 also broadcasts the radio channels Radio 1, Radio 2, 3 FM, Radio 4, Radio 5, Radio 6, Concertzender, FunX and also a regional channel. As of June 2011, the Dutch DVB-T service had 29 TV channels and 20 radio channels (including free to air channels). DVB-T2 will be introduced during 2019/2020.
Norway
In Norway, DVB-T broadcasting is marketed under RiksTV (encrypted pay channels) and NRK (unencrypted public channels). DVB-T broadcasting via the terrestrial network began in November 2007, and has subsequently been rolled out one part of the country at a time. The Norwegian implementation of DVB-T is different from most others, as it uses H.264 with HE-AAC audio encoding, while most other countries have adapted the less recent MPEG-2 standard. Notably most DVB software for PC has problems with this, though in late 2007 compatible software was released, like DVBViewer using the libfaad2 library. Sony has released several HDTVs (Bravia W3000, X3000, X3500, E4000, V4500, W4000, W4500, X4500) that support Norway's DVB-T implementation without use of a separate set-top box, and Sagem ITD91 HD, Grundig DTR 8720 STBs are others.
Poland
Currently, Poland uses the DVB-T2 standard with HEVC encoding. Analogue broadcast switch-off started on 7 November 2012 and was completed on 23 July 2013.
Portugal
Portugal follows the DVB-T implementation, using H.264 with AAC audio encoding. It has been live since 29 April 2009 and the switch-off date for all analog signals was on 26 April 2012.
Romania
Romania started digital terrestrial broadcasting in 2005 but it was virtually unknown by many people in Romania due to the lack of content, cable TV and satellite TV being far more popular, however it was the first platform to deliver HD content. Today, Romania is using DVB-T2 as terrestrial standard, but also DVB-S/S2, and DVB-C which is extremely popular. The only analogue broadcast remains on cable. Romania adopted the DVB-T2 standard in 2016 after a series of tests with mpeg2, mpeg4 on DVB-T, and has today fully implemented DVB-T2. DVB-C, which was introduced in late 2005, still remains with mpeg2 on SD content and mpeg4 on HD content. DVB-S (introduced in 2004 focus sat being the first such platform) is used in basic packages with standard definition content, while DVB-S2 set top boxes are provided for both SD and HD content.
Russia
Fully switched to digital in 2019, Russia uses the DVB-T2 standard for broadcasting 2 channel packs with about ten main national radio and TV channels (Channel One, Rossiya 1/2/K/24, NTV, Radio Mayak, Radio Rossii etc.
Spain
Quiero TV started digital terrestrial broadcasting in 2000 as pay television. The platform closed three years later after gaining 200,000 subscribers. The frequencies used by Quiero TV were used from 2005 to simulcast free-to-air analogue broadcast as DVB-T, under the name "TDT". The service started with 20 free-to-air national TV channels as well as numerous regional and local channels. Analogue broadcast ended on 2010 after getting 100% digital coverage. Some of the analogue frequencies were used to increase the number of channels and simulcast some of them in HD. Since February 14th, 2024, all channels will be required to broadcast exclusively in HD. Frequencies of SD channels will be used to simulcast some of them in 4K using DVB-T2.
United Kingdom
In the UK DVB-T has been adopted for broadcast of standard definition terrestrial programming, as well as a single DVB-T2 multiplex for high-definition programming. The UK terminated all analogue terrestrial broadcasts by the end of 2012. The vast majority of channels are available free-to-air through the Freeview service. DVB-T was also used for the now-defunct ONDigital/ITV Digital and Top Up TV service.
All satellite programming (some of which is available free-to-air via Freesat or free-to-view via Freesat from Sky; the remainder requires a subscription to Sky), is broadcast using either DVB-S or DVB-S2.
Subscription-based cable television from Virgin Media uses DVB-C.
North America
In North America, DVB-S is often used in encoding and video compression of digital satellite communications alongside Hughes DSS. Unlike Motorola's DigiCipher 2 standard, DVB has a wider adoption in terms of the number of manufacturers of receivers. Terrestrial digital television broadcasts in Canada, Mexico, El Salvador, Honduras, and the United States use ATSC encoding with 8VSB modulation instead of DVB-T with COFDM. Television newsgathering links from mobile vans to central receive points (often on mountaintops or tall buildings) use DVB-T with COFDM in the 2 GHz frequency band.
Oceania
Australia
In Australia, DVB broadcasting is marketed under the Freeview brand name, and more recently 'Freeview Plus', denoting the integration of online HbbTV and EPG in certain DVB devices. Regular broadcasts began in January 2001 using MPEG 2 video and MPEG 1 audio in SD and HD.
Changes to broadcasting rules have enabled broadcasters to offer multi-channeling, prompting broadcasters to use H.264 video with MPEG 1 or AAC audio encoding for some secondary channels.
Specifications for HD channels now differ depending on the broadcaster. ABC, Nine and Ten use 1920x1080i MPEG 4 video with Dolby Digital audio. Seven and SBS use 1440x1080i MPEG 2 video with Dolby Digital and MPEG 1 respectively.
New Zealand
In New Zealand, DVB broadcasting is marketed under the Freeview brand name. SD MPEG-2 DVB-S broadcasts via satellite began on 2 May 2007 and DVB-T (terrestrial) broadcasts began April 2008 broadcasting in HD H.264 video with HE-AAC audio.
South America
Colombia
Since 2008, Colombia has adopted as a public policy the decision to migrate from the analog television implemented in 1954 to Digital Terrestrial Television (DVB-T2). This measure allows the viewers access to the open television (OTA) of public and private channels, with video quality in HD. As planned, analogue television broadcasts will end in 2021.
DVB compliant products
Companies that manufacture a product which is compliant to one or more DVB standards have the option of registering a declaration of conformity for that product. Wherever the DVB trademark is used in relation to a product – be it a broadcast, a service, an application or equipment – the product must be registered with the DVB project office.
| Technology | Broadcasting | null |
146983 | https://en.wikipedia.org/wiki/Earth%27s%20magnetic%20field | Earth's magnetic field | Earth's magnetic field, also known as the geomagnetic field, is the magnetic field that extends from Earth's interior out into space, where it interacts with the solar wind, a stream of charged particles emanating from the Sun. The magnetic field is generated by electric currents due to the motion of convection currents of a mixture of molten iron and nickel in Earth's outer core: these convection currents are caused by heat escaping from the core, a natural process called a geodynamo.
The magnitude of Earth's magnetic field at its surface ranges from . As an approximation, it is represented by a field of a magnetic dipole currently tilted at an angle of about 11° with respect to Earth's rotational axis, as if there were an enormous bar magnet placed at that angle through the center of Earth. The North geomagnetic pole (Ellesmere Island, Nunavut, Canada) actually represents the South pole of Earth's magnetic field, and conversely the South geomagnetic pole corresponds to the north pole of Earth's magnetic field (because opposite magnetic poles attract and the north end of a magnet, like a compass needle, points toward Earth's South magnetic field.)
While the North and South magnetic poles are usually located near the geographic poles, they slowly and continuously move over geological time scales, but sufficiently slowly for ordinary compasses to remain useful for navigation. However, at irregular intervals averaging several hundred thousand years, Earth's field reverses and the North and South Magnetic Poles abruptly switch places. These reversals of the geomagnetic poles leave a record in rocks that are of value to paleomagnetists in calculating geomagnetic fields in the past. Such information in turn is helpful in studying the motions of continents and ocean floors. The magnetosphere is defined by the extent of Earth's magnetic field in space or geospace. It extends above the ionosphere, several tens of thousands of kilometres into space, protecting Earth from the charged particles of the solar wind and cosmic rays that would otherwise strip away the upper atmosphere, including the ozone layer that protects Earth from harmful ultraviolet radiation.
Significance
Earth's magnetic field deflects most of the solar wind, whose charged particles would otherwise strip away the ozone layer that protects the Earth from harmful ultraviolet radiation. One stripping mechanism is for gas to be caught in bubbles of the magnetic field, which are ripped off by solar winds. Calculations of the loss of carbon dioxide from the atmosphere of Mars, resulting from scavenging of ions by the solar wind, indicate that the dissipation of the magnetic field of Mars caused a near total loss of its atmosphere.
The study of the past magnetic field of the Earth is known as paleomagnetism. The polarity of the Earth's magnetic field is recorded in igneous rocks, and reversals of the field are thus detectable as "stripes" centered on mid-ocean ridges where the sea floor is spreading, while the stability of the geomagnetic poles between reversals has allowed paleomagnetism to track the past motion of continents. Reversals also provide the basis for magnetostratigraphy, a way of dating rocks and sediments. The field also magnetizes the crust, and magnetic anomalies can be used to search for deposits of metal ores.
Humans have used compasses for direction finding since the 11th century A.D. and for navigation since the 12th century. Although the magnetic declination does shift with time, this wandering is slow enough that a simple compass can remain useful for navigation. Using magnetoreception, various other organisms, ranging from some types of bacteria to pigeons, use the Earth's magnetic field for orientation and navigation.
Characteristics
At any location, the Earth's magnetic field can be represented by a three-dimensional vector. A typical procedure for measuring its direction is to use a compass to determine the direction of magnetic North. Its angle relative to true North is the declination () or variation. Facing magnetic North, the angle the field makes with the horizontal is the inclination () or magnetic dip. The intensity () of the field is proportional to the force it exerts on a magnet. Another common representation is in (North), (East) and (Down) coordinates.
Intensity
The intensity of the field is often measured in gauss (G), but is generally reported in microteslas (μT), with 1 G = 100 μT. A nanotesla is also referred to as a gamma (γ). The Earth's field ranges between approximately . By comparison, a strong refrigerator magnet has a field of about .
A map of intensity contours is called an isodynamic chart. As the World Magnetic Model shows, the intensity tends to decrease from the poles to the equator. A minimum intensity occurs in the South Atlantic Anomaly over South America while there are maxima over northern Canada, Siberia, and the coast of Antarctica south of Australia.
The intensity of the magnetic field is subject to change over time. A 2021 paleomagnetic study from the University of Liverpool contributed to a growing body of evidence that the Earth's magnetic field cycles with intensity every 200 million years. The lead author stated that "Our findings, when considered alongside the existing datasets, support the existence of an approximately 200-million-year-long cycle in the strength of the Earth's magnetic field related to deep Earth processes."
Inclination
The inclination is given by an angle that can assume values between −90° (up) to 90° (down). In the northern hemisphere, the field points downwards. It is straight down at the North Magnetic Pole and rotates upwards as the latitude decreases until it is horizontal (0°) at the magnetic equator. It continues to rotate upwards until it is straight up at the South Magnetic Pole. Inclination can be measured with a dip circle.
An isoclinic chart (map of inclination contours) for the Earth's magnetic field is shown below.
Declination
Declination is positive for an eastward deviation of the field relative to true north. It can be estimated by comparing the magnetic north–south heading on a compass with the direction of a celestial pole. Maps typically include information on the declination as an angle or a small diagram showing the relationship between magnetic north and true north. Information on declination for a region can be represented by a chart with isogonic lines (contour lines with each line representing a fixed declination).
Geographical variation
Dipolar approximation
Near the surface of the Earth, its magnetic field can be closely approximated by the field of a magnetic dipole positioned at the center of the Earth and tilted at an angle of about 11° with respect to the rotational axis of the Earth. The dipole is roughly equivalent to a powerful bar magnet, with its south pole pointing towards the geomagnetic North Pole. This may seem surprising, but the north pole of a magnet is so defined because, if allowed to rotate freely, it points roughly northward (in the geographic sense). Since the north pole of a magnet attracts the south poles of other magnets and repels the north poles, it must be attracted to the south pole of Earth's magnet. The dipolar field accounts for 80–90% of the field in most locations.
Magnetic poles
Historically, the north and south poles of a magnet were first defined by the Earth's magnetic field, not vice versa, since one of the first uses for a magnet was as a compass needle. A magnet's North pole is defined as the pole that is attracted by the Earth's North Magnetic Pole, in the arctic region, when the magnet is suspended so it can turn freely. Since opposite poles attract, the North Magnetic Pole of the Earth is really the south pole of its magnetic field (the place where the field is directed downward into the Earth).
The positions of the magnetic poles can be defined in at least two ways: locally or globally. The local definition is the point where the magnetic field is vertical. This can be determined by measuring the inclination. The inclination of the Earth's field is 90° (downwards) at the North Magnetic Pole and –90° (upwards) at the South Magnetic Pole. The two poles wander independently of each other and are not directly opposite each other on the globe. Movements of up to per year have been observed for the North Magnetic Pole. Over the last 180 years, the North Magnetic Pole has been migrating northwestward, from Cape Adelaide in the Boothia Peninsula in 1831 to from Resolute Bay in 2001. The magnetic equator is the line where the inclination is zero (the magnetic field is horizontal).
The global definition of the Earth's field is based on a mathematical model. If a line is drawn through the center of the Earth, parallel to the moment of the best-fitting magnetic dipole, the two positions where it intersects the Earth's surface are called the North and South geomagnetic poles. If the Earth's magnetic field were perfectly dipolar, the geomagnetic poles and magnetic dip poles would coincide and compasses would point towards them. However, the Earth's field has a significant non-dipolar contribution, so the poles do not coincide and compasses do not generally point at either.
Magnetosphere
Earth's magnetic field, predominantly dipolar at its surface, is distorted further out by the solar wind. This is a stream of charged particles leaving the Sun's corona and accelerating to a speed of 200 to 1000 kilometres per second. They carry with them a magnetic field, the interplanetary magnetic field (IMF).
The solar wind exerts a pressure, and if it could reach Earth's atmosphere it would erode it. However, it is kept away by the pressure of the Earth's magnetic field. The magnetopause, the area where the pressures balance, is the boundary of the magnetosphere. Despite its name, the magnetosphere is asymmetric, with the sunward side being about 10 Earth radii out but with the other side stretching out in a magnetotail that extends beyond 200 Earth radii.
Sunward of the magnetopause is the bow shock, the area where the solar wind slows abruptly.
Inside the magnetosphere is the plasmasphere, a donut-shaped region containing low-energy charged particles, or plasma. This region begins at a height of 60 km, extends up to 3 or 4 Earth radii, and includes the ionosphere. This region rotates with the Earth. There are also two concentric tire-shaped regions, called the Van Allen radiation belts, with high-energy ions (energies from 0.1 to 10 MeV). The inner belt is 1–2 Earth radii out while the outer belt is at 4–7 Earth radii. The plasmasphere and Van Allen belts have partial overlap, with the extent of overlap varying greatly with solar activity.
As well as deflecting the solar wind, the Earth's magnetic field deflects cosmic rays, high-energy charged particles that are mostly from outside the Solar System. Many cosmic rays are kept out of the Solar System by the Sun's magnetosphere, or heliosphere. By contrast, astronauts on the Moon risk exposure to radiation. Anyone who had been on the Moon's surface during a particularly violent solar eruption in 2005 would have received a lethal dose.
Some of the charged particles do get into the magnetosphere. These spiral around field lines, bouncing back and forth between the poles several times per second. In addition, positive ions slowly drift westward and negative ions drift eastward, giving rise to a ring current. This current reduces the magnetic field at the Earth's surface. Particles that penetrate the ionosphere and collide with the atoms there give rise to the lights of the aurorae while also emitting X-rays.
The varying conditions in the magnetosphere, known as space weather, are largely driven by solar activity. If the solar wind is weak, the magnetosphere expands; while if it is strong, it compresses the magnetosphere and more of it gets in. Periods of particularly intense activity, called geomagnetic storms, can occur when a coronal mass ejection erupts above the Sun and sends a shock wave through the Solar System. Such a wave can take just two days to reach the Earth. Geomagnetic storms can cause a lot of disruption; the "Halloween" storm of 2003 damaged more than a third of NASA's satellites. The largest documented storm, the Carrington Event, occurred in 1859. It induced currents strong enough to disrupt telegraph lines, and aurorae were reported as far south as Hawaii.
Time dependence
Short-term variations
The geomagnetic field changes on time scales from milliseconds to millions of years. Shorter time scales mostly arise from currents in the ionosphere (ionospheric dynamo region) and magnetosphere, and some changes can be traced to geomagnetic storms or daily variations in currents. Changes over time scales of a year or more mostly reflect changes in the Earth's interior, particularly the iron-rich core.
Frequently, the Earth's magnetosphere is hit by solar flares causing geomagnetic storms, provoking displays of aurorae. The short-term instability of the magnetic field is measured with the K-index.
Data from THEMIS show that the magnetic field, which interacts with the solar wind, is reduced when the magnetic orientation is aligned between Sun and Earth – opposite to the previous hypothesis. During forthcoming solar storms, this could result in blackouts and disruptions in artificial satellites.
Secular variation
Changes in Earth's magnetic field on a time scale of a year or more are referred to as secular variation. Over hundreds of years, magnetic declination is observed to vary over tens of degrees. The animation shows how global declinations have changed over the last few centuries.
The direction and intensity of the dipole change over time. Over the last two centuries the dipole strength has been decreasing at a rate of about 6.3% per century. At this rate of decrease, the field would be negligible in about 1600 years. However, this strength is about average for the last 7 thousand years, and the current rate of change is not unusual.
A prominent feature in the non-dipolar part of the secular variation is a westward drift at a rate of about 0.2° per year. This drift is not the same everywhere and has varied over time. The globally averaged drift has been westward since about 1400 AD but eastward between about 1000 AD and 1400 AD.
Changes that predate magnetic observatories are recorded in archaeological and geological materials. Such changes are referred to as paleomagnetic secular variation or paleosecular variation (PSV). The records typically include long periods of small change with occasional large changes reflecting geomagnetic excursions and reversals.
A 1995 study of lava flows on Steens Mountain, Oregon appeared to suggest the magnetic field once shifted at a rate of up to 6° per day at some time in Earth's history, a surprising result. However, in 2014 one of the original authors published a new study which found the results were actually due to the continuous thermal demagnitization of the lava, not to a shift in the magnetic field.
In July 2020 scientists report that analysis of simulations and a recent observational field model show that maximum rates of directional change of Earth's magnetic field reached ~10° per year – almost 100 times faster than current changes and 10 times faster than previously thought.
Magnetic field reversals
Although generally Earth's field is approximately dipolar, with an axis that is nearly aligned with the rotational axis, occasionally the North and South geomagnetic poles trade places. Evidence for these geomagnetic reversals can be found in basalts, sediment cores taken from the ocean floors, and seafloor magnetic anomalies. Reversals occur nearly randomly in time, with intervals between reversals ranging from less than 0.1 million years to as much as 50 million years. The most recent geomagnetic reversal, called the Brunhes–Matuyama reversal, occurred about 780,000 years ago. A related phenomenon, a geomagnetic excursion, takes the dipole axis across the equator and then back to the original polarity. The Laschamp event is an example of an excursion, occurring during the last ice age (41,000 years ago).
The past magnetic field is recorded mostly by strongly magnetic minerals, particularly iron oxides such as magnetite, that can carry a permanent magnetic moment. This remanent magnetization, or remanence, can be acquired in more than one way. In lava flows, the direction of the field is "frozen" in small minerals as they cool, giving rise to a thermoremanent magnetization. In sediments, the orientation of magnetic particles acquires a slight bias towards the magnetic field as they are deposited on an ocean floor or lake bottom. This is called detrital remanent magnetization.
Thermoremanent magnetization is the main source of the magnetic anomalies around mid-ocean ridges. As the seafloor spreads, magma wells up from the mantle, cools to form new basaltic crust on both sides of the ridge, and is carried away from it by seafloor spreading. As it cools, it records the direction of the Earth's field. When the Earth's field reverses, new basalt records the reversed direction. The result is a series of stripes that are symmetric about the ridge. A ship towing a magnetometer on the surface of the ocean can detect these stripes and infer the age of the ocean floor below. This provides information on the rate at which seafloor has spread in the past.
Radiometric dating of lava flows has been used to establish a geomagnetic polarity time scale, part of which is shown in the image. This forms the basis of magnetostratigraphy, a geophysical correlation technique that can be used to date both sedimentary and volcanic sequences as well as the seafloor magnetic anomalies.
Earliest appearance
Paleomagnetic studies of Paleoarchean lava in Australia and conglomerate in South Africa have concluded that the magnetic field has been present since at least about . In 2024 researchers published evidence from Greenland for the existence of the magnetic field as early as 3,700 million years ago.
Future
Starting in the late 1800s and throughout the 1900s and later, the overall geomagnetic field has become weaker; the present strong deterioration corresponds to a 10–15% decline and has accelerated since 2000; geomagnetic intensity has declined almost continuously from a maximum 35% above the modern value, from circa year 1 AD. The rate of decrease and the current strength are within the normal range of variation, as shown by the record of past magnetic fields recorded in rocks.
The nature of Earth's magnetic field is one of heteroscedastic (seemingly random) fluctuation. An instantaneous measurement of it, or several measurements of it across the span of decades or centuries, are not sufficient to extrapolate an overall trend in the field strength. It has gone up and down in the past for unknown reasons. Also, noting the local intensity of the dipole field (or its fluctuation) is insufficient to characterize Earth's magnetic field as a whole, as it is not strictly a dipole field. The dipole component of Earth's field can diminish even while the total magnetic field remains the same or increases.
The Earth's magnetic north pole is drifting from northern Canada towards Siberia with a presently accelerating rate— per year at the beginning of the 1900s, up to per year in 2003, and since then has only accelerated.
Physical origin
Earth's core and the geodynamo
The Earth's magnetic field is believed to be generated by electric currents in the conductive iron alloys of its core, created by convection currents due to heat escaping from the core.
The Earth and most of the planets in the Solar System, as well as the Sun and other stars, all generate magnetic fields through the motion of electrically conducting fluids. The Earth's field originates in its core. This is a region of iron alloys extending to about 3400 km (the radius of the Earth is 6370 km). It is divided into a solid inner core, with a radius of 1220 km, and a liquid outer core. The motion of the liquid in the outer core is driven by heat flow from the inner core, which is about , to the core-mantle boundary, which is about . The heat is generated by potential energy released by heavier materials sinking toward the core (planetary differentiation, the iron catastrophe) as well as decay of radioactive elements in the interior. The pattern of flow is organized by the rotation of the Earth and the presence of the solid inner core.
The mechanism by which the Earth generates a magnetic field is known as a geodynamo. The magnetic field is generated by a feedback loop: current loops generate magnetic fields (Ampère's circuital law); a changing magnetic field generates an electric field (Faraday's law); and the electric and magnetic fields exert a force on the charges that are flowing in currents (the Lorentz force). These effects can be combined in a partial differential equation for the magnetic field called the magnetic induction equation,
where is the velocity of the fluid; is the magnetic B-field; and is the magnetic diffusivity, which is the reciprocal of the product of the electrical conductivity and the permeability . The term is the partial derivative of the field with respect to time; is the Laplace operator, is the curl operator, and is the vector product.
The first term on the right hand side of the induction equation is a diffusion term. In a stationary fluid, the magnetic field declines and any concentrations of field spread out. If the Earth's dynamo shut off, the dipole part would disappear in a few tens of thousands of years.
In a perfect conductor (), there would be no diffusion. By Lenz's law, any change in the magnetic field would be immediately opposed by currents, so the flux through a given volume of fluid could not change. As the fluid moved, the magnetic field would go with it. The theorem describing this effect is called the frozen-in-field theorem. Even in a fluid with a finite conductivity, new field is generated by stretching field lines as the fluid moves in ways that deform it. This process could go on generating new field indefinitely, were it not that as the magnetic field increases in strength, it resists fluid motion.
The motion of the fluid is sustained by convection, motion driven by buoyancy. The temperature increases towards the center of the Earth, and the higher temperature of the fluid lower down makes it buoyant. This buoyancy is enhanced by chemical separation: As the core cools, some of the molten iron solidifies and is plated to the inner core. In the process, lighter elements are left behind in the fluid, making it lighter. This is called compositional convection. A Coriolis effect, caused by the overall planetary rotation, tends to organize the flow into rolls aligned along the north–south polar axis.
A dynamo can amplify a magnetic field, but it needs a "seed" field to get it started. For the Earth, this could have been an external magnetic field. Early in its history the Sun went through a T-Tauri phase in which the solar wind would have had a magnetic field orders of magnitude larger than the present solar wind. However, much of the field may have been screened out by the Earth's mantle. An alternative source is currents in the core-mantle boundary driven by chemical reactions or variations in thermal or electric conductivity. Such effects may still provide a small bias that are part of the boundary conditions for the geodynamo.
The average magnetic field in the Earth's outer core was calculated to be 25 gauss, 50 times stronger than the field at the surface.
Numerical models
Simulating the geodynamo by computer requires numerically solving a set of nonlinear partial differential equations for the magnetohydrodynamics (MHD) of the Earth's interior. Simulation of the MHD equations is performed on a 3D grid of points and the fineness of the grid, which in part determines the realism of the solutions, is limited mainly by computer power. For decades, theorists were confined to creating kinematic dynamo computer models in which the fluid motion is chosen in advance and the effect on the magnetic field calculated. Kinematic dynamo theory was mainly a matter of trying different flow geometries and testing whether such geometries could sustain a dynamo.
The first self-consistent dynamo models, ones that determine both the fluid motions and the magnetic field, were developed by two groups in 1995, one in Japan and one in the United States. The latter received attention because it successfully reproduced some of the characteristics of the Earth's field, including geomagnetic reversals.
Effect of ocean tides
The oceans contribute to Earth's magnetic field. Seawater is an electrical conductor, and therefore interacts with the magnetic field. As the tides cycle around the ocean basins, the ocean water essentially tries to pull the geomagnetic field lines along. Because the salty water is only slightly conductive, the interaction is relatively weak: the strongest component is from the regular lunar tide that happens about twice per day (M2). Other contributions come from ocean swell, eddies, and even tsunamis.
The strength of the interaction depends also on the temperature of the ocean water. The entire heat stored in the ocean can now be inferred from observations of the Earth's magnetic field.
Currents in the ionosphere and magnetosphere
Electric currents induced in the ionosphere generate magnetic fields (ionospheric dynamo region). Such a field is always generated near where the atmosphere is closest to the Sun, causing daily alterations that can deflect surface magnetic fields by as much as 1°. Typical daily variations of field strength are about 25 nT (one part in 2000), with variations over a few seconds of typically around 1 nT (one part in 50,000).
Measurement and analysis
Detection
The Earth's magnetic field strength was measured by Carl Friedrich Gauss in 1832 and has been repeatedly measured since then, showing a relative decay of about 10% over the last 150 years. The Magsat satellite and later satellites have used 3-axis vector magnetometers to probe the 3-D structure of the Earth's magnetic field. The later Ørsted satellite allowed a comparison indicating a dynamic geodynamo in action that appears to be giving rise to an alternate pole under the Atlantic Ocean west of South Africa.
Governments sometimes operate units that specialize in measurement of the Earth's magnetic field. These are geomagnetic observatories, typically part of a national Geological survey, for example, the British Geological Survey's Eskdalemuir Observatory. Such observatories can measure and forecast magnetic conditions such as magnetic storms that sometimes affect communications, electric power, and other human activities.
The International Real-time Magnetic Observatory Network, with over 100 interlinked geomagnetic observatories around the world, has been recording the Earth's magnetic field since 1991.
The military determines local geomagnetic field characteristics, in order to detect anomalies in the natural background that might be caused by a significant metallic object such as a submerged submarine. Typically, these magnetic anomaly detectors are flown in aircraft like the UK's Nimrod or towed as an instrument or an array of instruments from surface ships.
Commercially, geophysical prospecting companies also use magnetic detectors to identify naturally occurring anomalies from ore bodies, such as the Kursk Magnetic Anomaly.
Crustal magnetic anomalies
Magnetometers detect minute deviations in the Earth's magnetic field caused by iron artifacts, kilns, some types of stone structures, and even ditches and middens in archaeological geophysics. Using magnetic instruments adapted from airborne magnetic anomaly detectors developed during World War II to detect submarines, the magnetic variations across the ocean floor have been mapped. Basalt — the iron-rich, volcanic rock making up the ocean floor — contains a strongly magnetic mineral (magnetite) and can locally distort compass readings. The distortion was recognized by Icelandic mariners as early as the late 18th century. More important, because the presence of magnetite gives the basalt measurable magnetic properties, these magnetic variations have provided another means to study the deep ocean floor. When newly formed rock cools, such magnetic materials record the Earth's magnetic field.
Statistical models
Each measurement of the magnetic field is at a particular place and time. If an accurate estimate of the field at some other place and time is needed, the measurements must be converted to a model and the model used to make predictions.
Spherical harmonics
The most common way of analyzing the global variations in the Earth's magnetic field is to fit the measurements to a set of spherical harmonics. This was first done by Carl Friedrich Gauss. Spherical harmonics are functions that oscillate over the surface of a sphere. They are the product of two functions, one that depends on latitude and one on longitude. The function of longitude is zero along zero or more great circles passing through the North and South Poles; the number of such nodal lines is the absolute value of the order . The function of latitude is zero along zero or more latitude circles; this plus the order is equal to the degree ℓ. Each harmonic is equivalent to a particular arrangement of magnetic charges at the center of the Earth. A monopole is an isolated magnetic charge, which has never been observed. A dipole is equivalent to two opposing charges brought close together and a quadrupole to two dipoles brought together. A quadrupole field is shown in the lower figure on the right.
Spherical harmonics can represent any scalar field (function of position) that satisfies certain properties. A magnetic field is a vector field, but if it is expressed in Cartesian components , each component is the derivative of the same scalar function called the magnetic potential. Analyses of the Earth's magnetic field use a modified version of the usual spherical harmonics that differ by a multiplicative factor. A least-squares fit to the magnetic field measurements gives the Earth's field as the sum of spherical harmonics, each multiplied by the best-fitting Gauss coefficient or .
The lowest-degree Gauss coefficient, , gives the contribution of an isolated magnetic charge, so it is zero. The next three coefficients – , , and – determine the direction and magnitude of the dipole contribution. The best fitting dipole is tilted at an angle of about 10° with respect to the rotational axis, as described earlier.
Radial dependence
Spherical harmonic analysis can be used to distinguish internal from external sources if measurements are available at more than one height (for example, ground observatories and satellites). In that case, each term with coefficient or can be split into two terms: one that decreases with radius as and one that increases with radius as . The increasing terms fit the external sources (currents in the ionosphere and magnetosphere). However, averaged over a few years the external contributions average to zero.
The remaining terms predict that the potential of a dipole source () drops off as . The magnetic field, being a derivative of the potential, drops off as . Quadrupole terms drop off as , and higher order terms drop off increasingly rapidly with the radius. The radius of the outer core is about half of the radius of the Earth. If the field at the core-mantle boundary is fit to spherical harmonics, the dipole part is smaller by a factor of about 8 at the surface, the quadrupole part by a factor of 16, and so on. Thus, only the components with large wavelengths can be noticeable at the surface. From a variety of arguments, it is usually assumed that only terms up to degree or less have their origin in the core. These have wavelengths of about or less. Smaller features are attributed to crustal anomalies.
Global models
The International Association of Geomagnetism and Aeronomy maintains a standard global field model called the International Geomagnetic Reference Field (IGRF). It is updated every five years. The 11th-generation model, IGRF11, was developed using data from satellites (Ørsted, CHAMP and SAC-C) and a world network of geomagnetic observatories. The spherical harmonic expansion was truncated at degree 10, with 120 coefficients, until 2000. Subsequent models are truncated at degree 13 (195 coefficients).
Another global field model, called the World Magnetic Model, is produced jointly by the United States National Centers for Environmental Information (formerly the National Geophysical Data Center) and the British Geological Survey. This model truncates at degree 12 (168 coefficients) with an approximate spatial resolution of 3,000 kilometers. It is the model used by the United States Department of Defense, the Ministry of Defence (United Kingdom), the United States Federal Aviation Administration (FAA), the North Atlantic Treaty Organization (NATO), and the International Hydrographic Organization as well as in many civilian navigation systems.
The above models only take into account the "main field" at the core-mantle boundary. Although generally good enough for navigation, higher-accuracy use cases require smaller-scale magnetic anomalies and other variations to be considered. Some examples are (see geomag.us ref for more):
The "comprehensive modeling" (CM) approach by the Goddard Space Flight Center (NASA and GSFC) and the Danish Space Research Institute. CM attempts to reconcile data with greatly varying temporal and spatial resolution from ground and satellite sources. The latest version as of 2022 is CM5 of 2016. It provides separate components for main field plus lithosphere (crustal), M2 tidal, and primary/induced magnetosphere/ionosphere variations.
The US National Centers for Environmental Information developed the Enhanced Magnetic Model (EMM), which extends to degree and order 790 and resolves magnetic anomalies down to a wavelength of 56 kilometers. It was compiled from satellite, marine, aeromagnetic and ground magnetic surveys. , the latest version, EMM2017, includes data from The European Space Agency's Swarm satellite mission.
For historical data about the main field, the IGRF may be used back to year 1900. A specialized GUFM1 model estimates back to year 1590 using ship's logs. Paleomagnetic research has produced models dating back to 10,000 BCE.
Biomagnetism
Animals, including birds and turtles, can detect the Earth's magnetic field, and use the field to navigate during migration. Some researchers have found that cows and wild deer tend to align their bodies north–south while relaxing, but not when the animals are under high-voltage power lines, suggesting that magnetism is responsible. Other researchers reported in 2011 that they could not replicate those findings using different Google Earth images.
Very weak electromagnetic fields disrupt the magnetic compass used by European robins and other songbirds, which use the Earth's magnetic field to navigate. Neither power lines nor cellphone signals are to blame for the electromagnetic field effect on the birds; instead, the culprits have frequencies between 2 kHz and 5 MHz. These include AM radio signals and ordinary electronic equipment that might be found in businesses or private homes.
| Physical sciences | Geophysics | null |
147003 | https://en.wikipedia.org/wiki/Hysteresis | Hysteresis | Hysteresis is the dependence of the state of a system on its history. For example, a magnet may have more than one possible magnetic moment in a given magnetic field, depending on how the field changed in the past. Plots of a single component of the moment often form a loop or hysteresis curve, where there are different values of one variable depending on the direction of change of another variable. This history dependence is the basis of memory in a hard disk drive and the remanence that retains a record of the Earth's magnetic field magnitude in the past. Hysteresis occurs in ferromagnetic and ferroelectric materials, as well as in the deformation of rubber bands and shape-memory alloys and many other natural phenomena. In natural systems, it is often associated with irreversible thermodynamic change such as phase transitions and with internal friction; and dissipation is a common side effect.
Hysteresis can be found in physics, chemistry, engineering, biology, and economics. It is incorporated in many artificial systems: for example, in thermostats and Schmitt triggers, it prevents unwanted frequent switching.
Hysteresis can be a dynamic lag between an input and an output that disappears if the input is varied more slowly; this is known as rate-dependent hysteresis. However, phenomena such as the magnetic hysteresis loops are mainly rate-independent, which makes a durable memory possible.
Systems with hysteresis are nonlinear, and can be mathematically challenging to model. Some hysteretic models, such as the Preisach model (originally applied to ferromagnetism) and the Bouc–Wen model, attempt to capture general features of hysteresis; and there are also phenomenological models for particular phenomena such as the Jiles–Atherton model for ferromagnetism.
It is difficult to define hysteresis precisely. Isaak D. Mayergoyz wrote "...the very meaning of hysteresis varies from one area to another, from paper to paper and from author to author. As a result, a stringent mathematical definition of hysteresis is needed in order to avoid confusion and ambiguity.".
Etymology and history
The term "hysteresis" is derived from , an Ancient Greek word meaning "deficiency" or "lagging behind". It was coined in 1881 by Sir James Alfred Ewing to describe the behaviour of magnetic materials.
Some early work on describing hysteresis in mechanical systems was performed by James Clerk Maxwell. Subsequently, hysteretic models have received significant attention in the works of Ferenc Preisach (Preisach model of hysteresis), Louis Néel and Douglas Hugh Everett in connection with magnetism and absorption. A more formal mathematical theory of systems with hysteresis was developed in the 1970s by a group of Russian mathematicians led by Mark Krasnosel'skii.
Types
Rate-dependent
One type of hysteresis is a lag between input and output. An example is a sinusoidal input that results in a sinusoidal output , but with a phase lag :
Such behavior can occur in linear systems, and a more general form of response is
where is the instantaneous response and is the impulse response to an impulse that occurred time units in the past. In the frequency domain, input and output are related by a complex generalized susceptibility that can be computed from ; it is mathematically equivalent to a transfer function in linear filter theory and analogue signal processing.
This kind of hysteresis is often referred to as rate-dependent hysteresis. If the input is reduced to zero, the output continues to respond for a finite time. This constitutes a memory of the past, but a limited one because it disappears as the output decays to zero. The phase lag depends on the frequency of the input, and goes to zero as the frequency decreases.
When rate-dependent hysteresis is due to dissipative effects like friction, it is associated with power loss.
Rate-independent
Systems with rate-independent hysteresis have a persistent memory of the past that remains after the transients have died out. The future development of such a system depends on the history of states visited, but does not fade as the events recede into the past. If an input variable cycles from to and back again, the output may be initially but a different value upon return. The values of depend on the path of values that passes through but not on the speed at which it traverses the path. Many authors restrict the term hysteresis to mean only rate-independent hysteresis. Hysteresis effects can be characterized using the Preisach model and the generalized Prandtl−Ishlinskii model.
In engineering
Control systems
In control systems, hysteresis can be used to filter signals so that the output reacts less rapidly than it otherwise would by taking recent system history into account. For example, a thermostat controlling a heater may switch the heater on when the temperature drops below A, but not turn it off until the temperature rises above B. (For instance, if one wishes to maintain a temperature of 20 °C then one might set the thermostat to turn the heater on when the temperature drops to below 18 °C and off when the temperature exceeds 22 °C).
Similarly, a pressure switch can be designed to exhibit hysteresis, with pressure set-points substituted for temperature thresholds.
Electronic circuits
Often, some amount of hysteresis is intentionally added to an electronic circuit to prevent unwanted rapid switching. This and similar techniques are used to compensate for contact bounce in switches, or noise in an electrical signal.
A Schmitt trigger is a simple electronic circuit that exhibits this property.
A latching relay uses a solenoid to actuate a ratcheting mechanism that keeps the relay closed even if power to the relay is terminated.
Some positive feedback from the output to one input of a comparator can increase the natural hysteresis (a function of its gain) it exhibits.
Hysteresis is essential to the workings of some memristors (circuit components which "remember" changes in the current passing through them by changing their resistance).
Hysteresis can be used when connecting arrays of elements such as nanoelectronics, electrochrome cells and memory effect devices using passive matrix addressing. Shortcuts are made between adjacent components (see crosstalk) and the hysteresis helps to keep the components in a particular state while the other components change states. Thus, all rows can be addressed at the same time instead of individually.
In the field of audio electronics, a noise gate often implements hysteresis intentionally to prevent the gate from "chattering" when signals close to its threshold are applied.
User interface design
A hysteresis is sometimes intentionally added to computer algorithms. The field of user interface design has borrowed the term hysteresis to refer to times when the state of the user interface intentionally lags behind the apparent user input. For example, a menu that was drawn in response to a mouse-over event may remain on-screen for a brief moment after the mouse has moved out of the trigger region and the menu region. This allows the user to move the mouse directly to an item on the menu, even if part of that direct mouse path is outside of both the trigger region and the menu region. For instance, right-clicking on the desktop in most Windows interfaces will create a menu that exhibits this behavior.
Aerodynamics
In aerodynamics, hysteresis can be observed when decreasing the angle of attack of a wing after stall, regarding the lift and drag coefficients. The angle of attack at which the flow on top of the wing reattaches is generally lower than the angle of attack at which the flow separates during the increase of the angle of attack.
Hydraulics
Hysteresis can be observed in the stage-flow relationship of a river during rapidly changing conditions such as passing of a flood wave. It is most pronounced in low gradient streams with steep leading edge hydrographs.
Backlash
Moving parts within machines, such as the components of a gear train, normally have a small gap between them, to allow movement and lubrication. As a consequence of this gap, any reversal in direction of a drive part will not be passed on immediately to the driven part. This unwanted delay is normally kept as small as practicable, and is usually called backlash. The amount of backlash will increase with time as the surfaces of moving parts wear.
In mechanics
Elastic hysteresis
In the elastic hysteresis of rubber, the area in the centre of a hysteresis loop is the energy dissipated due to material internal friction.
Elastic hysteresis was one of the first types of hysteresis to be examined.
The effect can be demonstrated using a rubber band with weights attached to it. If the top of a rubber band is hung on a hook and small weights are attached to the bottom of the band one at a time, it will stretch and get longer. As more weights are loaded onto it, the band will continue to stretch because the force the weights are exerting on the band is increasing. When each weight is taken off, or unloaded, the band will contract as the force is reduced. As the weights are taken off, each weight that produced a specific length as it was loaded onto the band now contracts less, resulting in a slightly longer length as it is unloaded. This is because the band does not obey Hooke's law perfectly. The hysteresis loop of an idealized rubber band is shown in the figure.
In terms of force, the rubber band was harder to stretch when it was being loaded than when it was being unloaded. In terms of time, when the band is unloaded, the effect (the length) lagged behind the cause (the force of the weights) because the length has not yet reached the value it had for the same weight during the loading part of the cycle. In terms of energy, more energy was required during the loading than the unloading, the excess energy being dissipated as thermal energy.
Elastic hysteresis is more pronounced when the loading and unloading is done quickly than when it is done slowly. Some materials such as hard metals don't show elastic hysteresis under a moderate load, whereas other hard materials like granite and marble do. Materials such as rubber exhibit a high degree of elastic hysteresis.
When the intrinsic hysteresis of rubber is being measured, the material can be considered to behave like a gas. When a rubber band is stretched it heats up, and if it is suddenly released, it cools down perceptibly. These effects correspond to a large hysteresis from the thermal exchange with the environment and a smaller hysteresis due to internal friction within the rubber. This proper, intrinsic hysteresis can be measured only if the rubber band is thermally isolated.
Small vehicle suspensions using rubber (or other elastomers) can achieve the dual function of springing and damping because rubber, unlike metal springs, has pronounced hysteresis and does not return all the absorbed compression energy on the rebound. Mountain bikes have made use of elastomer suspension, as did the original Mini car.
The primary cause of rolling resistance when a body (such as a ball, tire, or wheel) rolls on a surface is hysteresis. This is attributed to the viscoelastic characteristics of the material of the rolling body.
Contact angle hysteresis
The contact angle formed between a liquid and solid phase will exhibit a range of contact angles that are possible. There are two common methods for measuring this range of contact angles. The first method is referred to as the tilting base method. Once a drop is dispensed on the surface with the surface level, the surface is then tilted from 0° to 90°. As the drop is tilted, the downhill side will be in a state of imminent wetting while the uphill side will be in a state of imminent dewetting. As the tilt increases the downhill contact angle will increase and represents the advancing contact angle while the uphill side will decrease; this is the receding contact angle. The values for these angles just prior to the drop releasing will typically represent the advancing and receding contact angles. The difference between these two angles is the contact angle hysteresis.
The second method is often referred to as the add/remove volume method. When the maximum liquid volume is removed from the drop without the interfacial area decreasing the receding contact angle is thus measured. When volume is added to the maximum before the interfacial area increases, this is the advancing contact angle. As with the tilt method, the difference between the advancing and receding contact angles is the contact angle hysteresis. Most researchers prefer the tilt method; the add/remove method requires that a tip or needle stay embedded in the drop which can affect the accuracy of the values, especially the receding contact angle.
Bubble shape hysteresis
The equilibrium shapes of bubbles expanding and contracting on capillaries (blunt needles) can exhibit hysteresis depending on the relative magnitude of the maximum capillary pressure to ambient pressure, and the relative magnitude of the bubble volume at the maximum capillary pressure to the dead volume in the system. The bubble shape hysteresis is a consequence of gas compressibility, which causes the bubbles to behave differently across expansion and contraction. During expansion, bubbles undergo large non equilibrium jumps in volume, while during contraction the bubbles are more stable and undergo a relatively smaller jump in volume resulting in an asymmetry across expansion and contraction. The bubble shape hysteresis is qualitatively similar to the adsorption hysteresis, and as in the contact angle hysteresis, the interfacial properties play an important role in bubble shape hysteresis.
The existence of the bubble shape hysteresis has important consequences in interfacial rheology experiments involving bubbles. As a result of the hysteresis, not all sizes of the bubbles can be formed on a capillary. Further the gas compressibility causing the hysteresis leads to unintended complications in the phase relation between the applied changes in interfacial area to the expected interfacial stresses. These difficulties can be avoided by designing experimental systems to avoid the bubble shape hysteresis.
Adsorption hysteresis
Hysteresis can also occur during physical adsorption processes. In this type of hysteresis, the quantity adsorbed is different when gas is being added than it is when being removed. The specific causes of adsorption hysteresis are still an active area of research, but it is linked to differences in the nucleation and evaporation mechanisms inside mesopores. These mechanisms are further complicated by effects such as cavitation and pore blocking.
In physical adsorption, hysteresis is evidence of mesoporosity-indeed, the definition of mesopores (2–50 nm) is associated with the appearance (50 nm) and disappearance (2 nm) of mesoporosity in nitrogen adsorption isotherms as a function of Kelvin radius. An adsorption isotherm showing hysteresis is said to be of Type IV (for a wetting adsorbate) or Type V (for a non-wetting adsorbate), and hysteresis loops themselves are classified according to how symmetric the loop is. Adsorption hysteresis loops also have the unusual property that it is possible to scan within a hysteresis loop by reversing the direction of adsorption while on a point on the loop. The resulting scans are called "crossing", "converging", or "returning", depending on the shape of the isotherm at this point.
Matric potential hysteresis
The relationship between matric water potential and water content is the basis of the water retention curve. Matric potential measurements (Ψm) are converted to volumetric water content (θ) measurements based on a site or soil specific calibration curve. Hysteresis is a source of water content measurement error. Matric potential hysteresis arises from differences in wetting behaviour causing dry medium to re-wet; that is, it depends on the saturation history of the porous medium. Hysteretic behaviour means that, for example, at a matric potential (Ψm) of , the volumetric water content (θ) of a fine sandy soil matrix could be anything between 8% and 25%.
Tensiometers are directly influenced by this type of hysteresis. Two other types of sensors used to measure soil water matric potential are also influenced by hysteresis effects within the sensor itself. Resistance blocks, both nylon and gypsum based, measure matric potential as a function of electrical resistance. The relation between the sensor's electrical resistance and sensor matric potential is hysteretic. Thermocouples measure matric potential as a function of heat dissipation. Hysteresis occurs because measured heat dissipation depends on sensor water content, and the sensor water content–matric potential relationship is hysteretic. , only desorption curves are usually measured during calibration of soil moisture sensors. Despite the fact that it can be a source of significant error, the sensor specific effect of hysteresis is generally ignored.
In materials
Magnetic hysteresis
When an external magnetic field is applied to a ferromagnetic material such as iron, the atomic domains align themselves with it. Even when the field is removed, part of the alignment will be retained: the material has become magnetized. Once magnetized, the magnet will stay magnetized indefinitely. To demagnetize it requires heat or a magnetic field in the opposite direction. This is the effect that provides the element of memory in a hard disk drive.
The relationship between field strength and magnetization is not linear in such materials. If a magnet is demagnetized () and the relationship between and is plotted for increasing levels of field strength, follows the initial magnetization curve. This curve increases rapidly at first and then approaches an asymptote called magnetic saturation. If the magnetic field is now reduced monotonically, follows a different curve. At zero field strength, the magnetization is offset from the origin by an amount called the remanence. If the relationship is plotted for all strengths of applied magnetic field the result is a hysteresis loop called the main loop. The width of the middle section is twice the coercivity of the material.
A closer look at a magnetization curve generally reveals a series of small, random jumps in magnetization called Barkhausen jumps. This effect is due to crystallographic defects such as dislocations.
Magnetic hysteresis loops are not exclusive to materials with ferromagnetic ordering. Other magnetic orderings, such as spin glass ordering, also exhibit this phenomenon.
Physical origin
The phenomenon of hysteresis in ferromagnetic materials is the result of two effects: rotation of magnetization and changes in size or number of magnetic domains. In general, the magnetization varies (in direction but not magnitude) across a magnet, but in sufficiently small magnets, it does not. In these single-domain magnets, the magnetization responds to a magnetic field by rotating. Single-domain magnets are used wherever a strong, stable magnetization is needed (for example, magnetic recording).
Larger magnets are divided into regions called domains. Across each domain, the magnetization does not vary; but between domains are relatively thin domain walls in which the direction of magnetization rotates from the direction of one domain to another. If the magnetic field changes, the walls move, changing the relative sizes of the domains. Because the domains are not magnetized in the same direction, the magnetic moment per unit volume is smaller than it would be in a single-domain magnet; but domain walls involve rotation of only a small part of the magnetization, so it is much easier to change the magnetic moment. The magnetization can also change by addition or subtraction of domains (called nucleation and denucleation).
Magnetic hysteresis models
The most known empirical models in hysteresis are Preisach and Jiles-Atherton models. These models allow an accurate modeling of the hysteresis loop and are widely used in the industry. However, these models lose the connection with thermodynamics and the energy consistency is not ensured. A more recent model, with a more consistent thermodynamical foundation, is the vectorial incremental nonconservative consistent hysteresis (VINCH) model of Lavet et al. (2011)
Applications
There are a great variety of applications of the hysteresis in ferromagnets. Many of these make use of their ability to retain a memory, for example magnetic tape, hard disks, and credit cards. In these applications, hard magnets (high coercivity) like iron are desirable, such that as much energy is absorbed as possible during the write operation and the resultant magnetized information is not easily erased.
On the other hand, magnetically soft (low coercivity) iron is used for the cores in electromagnets. The low coercivity minimizes the energy loss associated with hysteresis, as the magnetic field periodically reverses in the presence of an alternating current. The low energy loss during a hysteresis loop is the reason why soft iron is used for transformer cores and electric motors.
Electrical hysteresis
Electrical hysteresis typically occurs in ferroelectric material, where domains of polarization contribute to the total polarization. Polarization is the electrical dipole moment (either C·m−2 or C·m). The mechanism, an organization of the polarization into domains, is similar to that of magnetic hysteresis.
Liquid–solid-phase transitions
Hysteresis manifests itself in state transitions when melting temperature and freezing temperature do not agree. For example, agar melts at and solidifies from . This is to say that once agar is melted at 85 °C, it retains a liquid state until cooled to 40 °C. Therefore, from the temperatures of 40 to 85 °C, agar can be either solid or liquid, depending on which state it was before.
In biology
Cell biology and genetics
Hysteresis in cell biology often follows bistable systems where the same input state can lead to two different, stable outputs. Where bistability can lead to digital, switch-like outputs from the continuous inputs of chemical concentrations and activities, hysteresis makes these systems more resistant to noise. These systems are often characterized by higher values of the input required to switch into a particular state as compared to the input required to stay in the state, allowing for a transition that is not continuously reversible, and thus less susceptible to noise.
Irreversible Hysteresis
In the case of mitosis, irreversibility is essential to maintain the overall integrity of the system such that we have three designated checkpoints to account for this: G1/S, G2/M, and the spindle checkpoint. Irreversible hysteresis in this context ensures that once a cell commits to a specific phase (e.g., entering mitosis or DNA replication), it does not revert to a previous phase, even if conditions or regulatory signals change. Based on the irreversible hysteresis curve, there does exit a [input] at which the cell jumps to the next bistable region, but there isn’t a [input] that allows the cell to revert to its previous bistable region. The cell has no way of reverting back down to its previous bistable region even when [input] is 0, demonstrating irreversibility. Positive feedback is critical for generating hysteresis in the cell cycle. For example: In the G2/M transition, active CDK1 promotes the activation of more CDK1 molecules by inhibiting Wee1 (an inhibitor) and activating Cdc25 (a phosphatase that activates CDK1). These loops lock the cell into its current state and amplifying the activation of CDK1. Positive feedback also serves to create a bistable system where CDK1 is either fully inactivated or fully activated. Hysteresis prevents the cell from oscillating between these two states from small perturbations in signal (input).
Reversible Hysteresis
A biochemical system that is under the control of reversible hysteresis has both forward and reverse trajectories. The system generally requires a higher [input] to proceed forward into the next bistable state then to exit from that stage. For example, cells undergoing cell division exhibit reversible hysteresis in that it takes a higher concentration of cyclins to switch them from G2 phase into mitosis than to stay in mitosis once begun
. Additionally, because the [cyclin] required to reverse the cell back to the G2 phase is much lower than the [cycilin] to enter mitosis, this improved the bistability of mitosis because it is more resistance to weak or transient signals. Small perturbations the [input] won’t be able to push the cell out of mitosis so easily.
History and Memory
In systems with bistability, the same input level can correspond to two distinct stable states (e.g., "low output" and "high output"). The actual state of the system depends on its history whether the input level was increasing (forward trajectory) or decreasing (backward trajectory). Thus, it is difficult to determine which state a cell is in if given only a bistability curve. The cell’s ability to "remember" its prior state ensures stability and prevents it from switching states unnecessarily due to minor fluctuations in input. This memory is often maintained through molecular feedback loops, such as positive feedback in signaling pathways, or the persistence of regulatory molecules like proteins or phosphorylated components. For example, the refractory period in action potentials is primarily controlled by history. Absolute refraction period prevents a volted-gated sodium channel from activating or refiring after it has just fired. This is because following the absolute refractory period, the neuron is less excitable due to hyperpolarization caused by potassium efflux. This molecular inhibitory feedback creates a memory for the neuron or cell, so that the neuron doesn’t fire too soon. As time passes, the neuron or cell will slowly lose the memory of having fired and will begin to fire again. Thus, memory is time-dependent, which is important in maintaining homeostasis and regulating many different biological processes.
Biochemical Systems: Regulating the Cell Cycle in Xenopus Laevis Egg Extracts
Cells advancing through the cell cycle must make an irreversible commitment to mitosis, ensuring they do not revert to interphase before successfully segregating their chromosomes. A mathematical model of cell-cycle progression in cell-free egg extracts from frogs suggests that hysteresis in the molecular control system drives these irreversible transitions into and out of mitosis. Here, Cdc2 (Cyclin-dependent kinase 1 or CDK1) is responsible for mitotic entry and exit such that binding of cyclin B forms a complex called Maturation-Promoting Factor (MPF). The activation threshold for mitotic entry was found to be between 32 and 40 nM cyclin B in the frog extracts while the inactivation threshold for exiting mitosis was lower, between 16 and 24 nM cyclin B. The higher threshold for mitotic entry compared to the lower threshold for mitotic exit indicates hysteresis, a hallmark of history-dependent behavior in the system. Concentrations between 24 and 32 nM cyclin B demonstrated bistability, where the system could exist in either interphase or mitosis, depending on its prior state (history). Though, the cell cycle isn’t completely irreversible, the difference in thresholds is enough for growth and survival of the cells.
Hysteric thresholds in biological systems are not definite and can be recalibrated. For example, unreplicated DNA or chromosomes inhibits Cdc25 phosphatase and maintains Wee1 kinase activity. This prevents the activation of Cyclin B-Cdc2, effectively raising the threshold for mitotic entry. As a result, the cell delays the transition to mitosis until replication is complete, ensuring genomic integrity. Other instances may be DNA damage and unattached chromosomes during the spindle assembly checkpoint.
Biochemical Systems: Regulating the Cell Cycle in Yeast
Biochemical systems can also show hysteresis-like output when slowly varying states that are not directly monitored are involved, as in the case of the cell cycle arrest in yeast exposed to mating pheromone. The proposed model is that α-factor, a yeast mating pheromone binds to its analog receptor on another yeast cell promoting transcription of Fus3 and promoting mating. Fus3 further promotes Far1 which inhibits Cln1/2, activators of the cell cycle. This is representative of a coherent feedforward loop that can modeled as a hysteresis curve.
Far1 transcription is the primary mechanism responsible for the hysteresis observed in cell-cycle reentry. The history of pheromone exposure influences the accumulation of Far1, which, in turn, determines the delay in cell-cycle reentry. Previous pulse experiments demonstrated that after exposure to high pheromone concentrations, cells enter a stabilized arrested state where reentry thresholds are elevated due to increased Far1-dependent inhibition of CDK activity. Even when pheromone levels drop to concentrations that would allow naive cells to reenter the cell cycle, pre-exposed cells take longer to resume proliferation. This delay reflects the history-dependent nature of hysteresis, where past exposure to high pheromone concentrations influences the current state. Hysteresis ensures that cells make robust and irreversible decisions about mating and proliferation in response to pheromone signals. It allows cells to "remember" high pheromone exposure, and this helps yeast cells adapt and stability their responses to environmental conditions, avoiding fast premature reentry into the cell cycle, the moment that pheromone signal dies down.
Additionally, the duration of cell cycle arrest depends not only on the final level of input Fus3, but also on the previously achieved Fus3 levels. This effect is achieved due to the slower time scales involved in the transcription of intermediate Far1, such that the total Far1 activity reaches its equilibrium value slowly, and for transient changes in Fus3 concentration, the response of the system depends on the Far1 concentration achieved with the transient value. Experiments in this type of hysteresis benefit from the ability to change the concentration of the inputs with time. The mechanisms are often elucidated by allowing independent control of the concentration of the key intermediate, for instance, by using an inducible promoter.
Biochemical systems can also show hysteresis-like output when slowly varying states that are not directly monitored are involved, as in the case of the cell cycle arrest in yeast exposed to mating pheromone. Here, the duration of cell cycle arrest depends not only on the final level of input Fus3, but also on the previously achieved Fus3 levels. This effect is achieved due to the slower time scales involved in the transcription of intermediate Far1, such that the total Far1 activity reaches its equilibrium value slowly, and for transient changes in Fus3 concentration, the response of the system depends on the Far1 concentration achieved with the transient value. Experiments in this type of hysteresis benefit from the ability to change the concentration of the inputs with time. The mechanisms are often elucidated by allowing independent control of the concentration of the key intermediate, for instance, by using an inducible promoter.
Darlington in his classic works on genetics discussed hysteresis of the chromosomes, by which he meant "failure of the external form of the chromosomes to respond immediately to the internal stresses due to changes in their molecular spiral", as they lie in a somewhat rigid medium in the limited space of the cell nucleus.
In developmental biology, cell type diversity is regulated by long range-acting signaling molecules called morphogens that pattern uniform pools of cells in a concentration- and time-dependent manner. The morphogen sonic hedgehog (Shh), for example, acts on limb bud and neural progenitors to induce expression of a set of homeodomain-containing transcription factors to subdivide these tissues into distinct domains. It has been shown that these tissues have a 'memory' of previous exposure to Shh.
In neural tissue, this hysteresis is regulated by a homeodomain (HD) feedback circuit that amplifies Shh signaling. In this circuit, expression of Gli transcription factors, the executors of the Shh pathway, is suppressed. Glis are processed to repressor forms (GliR) in the absence of Shh, but in the presence of Shh, a proportion of Glis are maintained as full-length proteins allowed to translocate to the nucleus, where they act as activators (GliA) of transcription. By reducing Gli expression then, the HD transcription factors reduce the total amount of Gli (GliT), so a higher proportion of GliT can be stabilized as GliA for the same concentration of Shh.
Immunology
There is some evidence that T cells exhibit hysteresis in that it takes a lower signal threshold to activate T cells that have been previously activated. Ras GTPase activation is required for downstream effector functions of activated T cells. Triggering of the T cell receptor induces high levels of Ras activation, which results in higher levels of GTP-bound (active) Ras at the cell surface. Since higher levels of active Ras have accumulated at the cell surface in T cells that have been previously stimulated by strong engagement of the T cell receptor, weaker subsequent T cell receptor signals received shortly afterwards will deliver the same level of activation due to the presence of higher levels of already activated Ras as compared to a naïve cell.
Neuroscience
The property by which some neurons do not return to their basal conditions from a stimulated condition immediately after removal of the stimulus is an example of hysteresis.
Neuropsychology
Neuropsychology, in exploring the neural correlates of consciousness, interfaces with neuroscience, although the complexity of the central nervous system is a challenge to its study (that is, its operation resists easy reduction). Context-dependent memory and state-dependent memory show hysteretic aspects of neurocognition.
Respiratory physiology
Lung hysteresis is evident when observing the compliance of a lung on inspiration versus expiration. The difference in compliance (Δvolume/Δpressure) is due to the additional energy required to overcome surface tension forces during inspiration to recruit and inflate additional alveoli.
The transpulmonary pressure vs Volume curve of inhalation is different from the Pressure vs Volume curve of exhalation, the difference being described as hysteresis. Lung volume at any given pressure during inhalation is less than the lung volume at any given pressure during exhalation.
Voice and speech physiology
A hysteresis effect may be observed in voicing onset versus offset. The threshold value of the subglottal pressure required to start the vocal fold vibration is lower than the threshold value at which the vibration stops, when other parameters are kept constant. In utterances of vowel-voiceless consonant-vowel sequences during speech, the intraoral pressure is lower at the voice onset of the second vowel compared to the voice offset of the first vowel, the oral airflow is lower, the transglottal pressure is larger and the glottal width is smaller.
Ecology and epidemiology
Hysteresis is a commonly encountered phenomenon in ecology and epidemiology, where the observed equilibrium of a system can not be predicted solely based on environmental variables, but also requires knowledge of the system's past history. Notable examples include the theory of spruce budworm outbreaks and behavioral-effects on disease transmission.
It is commonly examined in relation to critical transitions between ecosystem or community types in which dominant competitors or entire landscapes can change in a largely irreversible fashion.
In ocean and climate science
Complex ocean and climate models rely on the principle.
In economics
Economic systems can exhibit hysteresis. For example, export performance is subject to strong hysteresis effects: because of the fixed transportation costs it may take a big push to start a country's exports, but once the transition is made, not much may be required to keep them going.
When some negative shock reduces employment in a company or industry, fewer employed workers then remain. As usually the employed workers have the power to set wages, their reduced number incentivizes them to bargain for higher wages when the economy again gets better instead of letting the wage be at the equilibrium wage level, where the supply and demand of workers would match. This causes hysteresis: the unemployment becomes permanently higher after negative shocks.
Permanently higher unemployment
The idea of hysteresis is used extensively in the area of labor economics, specifically with reference to the unemployment rate. According to theories based on hysteresis, severe economic downturns (recession) and/or persistent stagnation (slow demand growth, usually after a recession) cause unemployed individuals to lose their job skills (commonly developed on the job) or to find that their skills have become obsolete, or become demotivated, disillusioned or depressed or lose job-seeking skills. In addition, employers may use time spent in unemployment as a screening tool, i.e., to weed out less desired employees in hiring decisions. Then, in times of an economic upturn, recovery, or "boom", the affected workers will not share in the prosperity, remaining unemployed for long periods (e.g., over 52 weeks). This makes unemployment "structural", i.e., extremely difficult to reduce simply by increasing the aggregate demand for products and labor without causing increased inflation. That is, it is possible that a ratchet effect in unemployment rates exists, so a short-term rise in unemployment rates tends to persist. For example, traditional anti-inflationary policy (the use of recession to fight inflation) leads to a permanently higher "natural" rate of unemployment (more scientifically known as the NAIRU). This occurs first because inflationary expectations are "sticky" downward due to wage and price rigidities (and so adapt slowly over time rather than being approximately correct as in theories of rational expectations) and second because labor markets do not clear instantly in response to unemployment.
The existence of hysteresis has been put forward as a possible explanation for the persistently high unemployment of many economies in the 1990s. Hysteresis has been invoked by Olivier Blanchard among others to explain the differences in long run unemployment rates between Europe and the United States. Labor market reform (usually meaning institutional change promoting more flexible wages, firing, and hiring) or strong demand-side economic growth may not therefore reduce this pool of long-term unemployed. Thus, specific targeted training programs are presented as a possible policy solution. However, the hysteresis hypothesis suggests such training programs are aided by persistently high demand for products (perhaps with incomes policies to avoid increased inflation), which reduces the transition costs out of unemployment and into paid employment easier.
Models
Hysteretic models are mathematical models capable of simulating complex nonlinear behavior (hysteresis) characterizing mechanical systems and materials used in different fields of engineering, such as aerospace, civil, and mechanical engineering. Some examples of mechanical systems and materials having hysteretic behavior are:
materials, such as steel, reinforced concrete, wood;
structural elements, such as steel, reinforced concrete, or wood joints;
devices, such as seismic isolators and dampers.
Each subject that involves hysteresis has models that are specific to the subject. In addition, there are hysteretic models that capture general features of many systems with hysteresis. An example is the Preisach model of hysteresis, which represents a hysteresis nonlinearity as a linear superposition of square loops called non-ideal relays. Many complex models of hysteresis arise from the simple parallel connection, or superposition, of elementary carriers of hysteresis termed hysterons.
A simple and intuitive parametric description of various hysteresis loops may be found in the Lapshin model. Along with the smooth loops, substitution of trapezoidal, triangular or rectangular pulses instead of the harmonic functions allows piecewise-linear hysteresis loops frequently used in discrete automatics to be built in the model. There are implementations of the hysteresis loop model in Mathcad and in R programming language.
The Bouc–Wen model of hysteresis is often used to describe non-linear hysteretic systems. It was introduced by Bouc and extended by Wen, who demonstrated its versatility by producing a variety of hysteretic patterns. This model is able to capture in analytical form, a range of shapes of hysteretic cycles which match the behaviour of a wide class of hysteretical systems; therefore, given its versability and mathematical tractability, the Bouc–Wen model has quickly gained popularity and has been extended and applied to a wide variety of engineering problems, including multi-degree-of-freedom (MDOF) systems, buildings, frames, bidirectional and torsional response of hysteretic systems two- and three-dimensional continua, and soil liquefaction among others. The Bouc–Wen model and its variants/extensions have been used in applications of structural control, in particular in the modeling of the behaviour of magnetorheological dampers, base isolation devices for buildings and other kinds of damping devices; it has also been used in the modelling and analysis of structures built of reinforced concrete, steel, masonry and timber.. The most important extension of Bouc-Wen Model was carried out by Baber and Noori and later by Noori and co-workers. That extended model, named, BWBN, can reproduce the complex shear pinching or slip-lock phenomenon that earlier model could not reproduce. The BWBN model has been widely used in a wide spectrum of applications and implementations are available in software such as OpenSees.
Hysteretic models may have a generalized displacement as input variable and a generalized force as output variable, or vice versa. In particular, in rate-independent hysteretic models, the output variable does not depend on the rate of variation of the input one.
Rate-independent hysteretic models can be classified into four different categories depending on the type of equation that needs to be solved to compute the output variable:
algebraic models
transcendental models
differential models
integral models
List of models
Some notable hysteretic models are listed below, along with their associated fields.
Bean's critical state model (magnetism)
Bouc–Wen model (structural engineering)
Ising model (magnetism)
Jiles–Atherton model (magnetism)
Novak–Tyson model (cell-cycle control)
Preisach model (magnetism)
Stoner–Wohlfarth model (magnetism)
Energy
When hysteresis occurs with extensive and intensive variables, the work done on the system is the area under the hysteresis graph.
| Physical sciences | Physics basics: General | Physics |
147019 | https://en.wikipedia.org/wiki/Seabird | Seabird | Seabirds (also known as marine birds) are birds that are adapted to life within the marine environment. While seabirds vary greatly in lifestyle, behaviour and physiology, they often exhibit striking convergent evolution, as the same environmental problems and feeding niches have resulted in similar adaptations. The first seabirds evolved in the Cretaceous period, while modern seabird families emerged in the Paleogene.
Seabirds generally live longer, breed later and have fewer young than other birds, but they invest a great deal of time in their young. Most species nest in colonies, varying in size from a few dozen birds to millions. Many species are famous for undertaking long annual migrations, crossing the equator or circumnavigating the Earth in some cases. They feed both at the ocean's surface and below it, and even on each other. Seabirds can be highly pelagic, coastal, or in some cases spend a part of the year away from the sea entirely.
Seabirds and humans have a long history together: They have provided food to hunters, guided fishermen to fishing stocks, and led sailors to land. Many species are currently threatened by human activities such as oil spills, nets, climate change and severe weather. Conservation efforts include the establishment of wildlife refuges and adjustments to fishing techniques.
Classification
There exists no single definition of which groups, families and species are seabirds, and most definitions are in some way arbitrary. Elizabeth Shreiber and Joanna Burger, two seabird scientists, said, "The one common characteristic that all seabirds share is that they feed in saltwater; but, as seems to be true with any statement in biology, some do not." However, by convention, all of the Sphenisciformes (penguins), all of the Phaethontiformes (tropicbirds), all of the Procellariiformes (albatrosses and petrels), all of the Suliformes (gannets, boobies, frigatebirds, and cormorants) except the darters, one family of the Pelecaniformes (pelicans), and some of the Charadriiformes (gulls, skuas, terns, auks, and skimmers) are classified as seabirds. The phalaropes are usually included as well, since although they are waders ("shorebirds" in North America), two of the three species (red and red-necked) are oceanic for nine months of the year, crossing the equator to feed pelagically.
Loons and grebes, which nest on lakes but winter at sea, are usually categorized as water birds, not seabirds. Although there are a number of sea ducks in the family Anatidae that are truly marine in the winter, by convention they are usually excluded from the seabird grouping. Many herons and waders (or shorebirds), such as crab-plovers, are also highly marine, living on the sea's edge (coast), but are also not treated as seabirds. Fish-eating birds of prey, such as sea eagles and ospreys, are also typically excluded, however tied to marine environments they may be. Some birds, such as darters and anhingas, are primarily found in freshwater habitats, but may occasionally venture into marine or coastal areas as well; such birds are generally not considered to be seabirds.
German ornithologist Gerald Mayr defined the "core waterbird" clade Aequornithes in 2010. This lineage gives rise to the Procellariiformes, Sphenisciformes, Suliformes, Pelecaniformes, Ciconiiformes (not seabirds), and Gaviiformes (not seabirds). The tropicbirds (Phaethontiformes) are part of the Eurypygimorphae lineage, which is sister to the Aequornithes; this clade also includes the non-seabird Eurypygiformes (kagu and sunbittern). The Charadriiformes are more distantly related to the other seabirds, being more closely related to the non-seabird Gruiformes (rails and cranes) and Opisthocomiformes (hoatzin) in the clade Gruae.
Evolution and fossil record
Seabirds, by virtue of living in a geologically depositional environment (that is, in the sea where sediments are readily laid down), are well represented in the fossil record. They are first known to occur in the Cretaceous period, the earliest being the Hesperornithes. These were flightless seabirds that could dive in a fashion similar to grebes and loons (using its feet to move underwater), but had beaks filled with sharp teeth. Other Cretaceous seabirds included the gull-like Ichthyornithes. Flying Cretaceous seabirds do not exceed wingspans of two meters; piscivorous pterosaurs occupied seagoing niches above this size.
While Hesperornis is not thought to have left descendants, the earliest modern seabirds also occurred in the Cretaceous, with a species called Tytthostonyx glauconiticus, which has features suggestive of Procellariiformes and Fregatidae. As a clade, the Aequornithes either became seabirds in a single transition in the Cretaceous or some lineages such as pelicans and frigatebirds adapted to sea living independently from freshwater-dwelling ancestors. In the Paleogene both pterosaurs and marine reptiles became extinct, allowing seabirds to expand ecologically. These post-extinction seas were dominated by early Procellariidae, giant penguins and two extinct families, the Pelagornithidae and the Plotopteridae (a group of large seabirds that looked like the penguins). Modern genera began their wide radiation in the Miocene, although the genus Puffinus (which includes today's Manx shearwater and sooty shearwater) might date back to the Oligocene. Within the Charadriiformes, the gulls and allies (Lari) became seabirds in the late Eocene, and then waders in the middle Miocene (Langhian). The highest diversity of seabirds apparently existed during the Late Miocene and the Pliocene. At the end of the latter, the oceanic food web had undergone a period of upheaval due to extinction of considerable numbers of marine species; subsequently, the spread of marine mammals seems to have prevented seabirds from reaching their erstwhile diversity.
Characteristics
Adaptations to life at sea
Seabirds have made numerous adaptations to living on and feeding in the sea. Wing morphology has been shaped by the niche an individual species or family has evolved, so that looking at a wing's shape and loading can tell a scientist about its life feeding behaviour. Longer wings and low wing loading are typical of more pelagic species, while diving species have shorter wings. Species such as the wandering albatross, which forage over huge areas of sea, have a reduced capacity for powered flight and are dependent on a type of gliding called dynamic soaring (where the wind deflected by waves provides lift) as well as slope soaring. Seabirds also almost always have webbed feet, to aid movement on the surface as well as assisting diving in some species. The Procellariiformes are unusual among birds in having a strong sense of smell, which is used to find widely distributed food in a vast ocean, and help distinguish familiar nest odours from unfamiliar ones.
Salt glands are used by seabirds to deal with the salt they ingest by drinking and feeding (particularly on crustaceans), and to help them osmoregulate. The excretions from these glands (which are positioned in the head of the birds, emerging from the nasal cavity) are almost pure sodium chloride.
With the exception of the cormorants and some terns, and in common with most other birds, all seabirds have waterproof plumage. However, compared to land birds, they have far more feathers protecting their bodies. This dense plumage is better able to protect the bird from getting wet, and cold is kept out by a dense layer of down feathers. The cormorants possess a layer of unique feathers that retain a smaller layer of air (compared to other diving birds) but otherwise soak up water. This allows them to swim without fighting the buoyancy that retaining air in the feathers causes, yet retain enough air to prevent the bird losing excessive heat through contact with water.
The plumage of most seabirds is less colourful than that of land birds, restricted in the main to variations of black, white or grey. A few species sport colourful plumes (such as the tropicbirds and some penguins), but most of the colour in seabirds appears in the bills and legs. The plumage of seabirds is thought in many cases to be for camouflage, both defensive (the colour of US Navy battleships is the same as that of Antarctic prions, and in both cases it reduces visibility at sea) and aggressive (the white underside possessed by many seabirds helps hide them from prey below). The usually black wing tips help prevent wear, as they contain melanins that help the feathers resist abrasion.
Diet and feeding
Seabirds evolved to exploit different food resources in the world's seas and oceans, and to a great extent, their physiology and behaviour have been shaped by their diet. These evolutionary forces have often caused species in different families and even orders to evolve similar strategies and adaptations to the same problems, leading to remarkable convergent evolution, such as that between auks and penguins. There are four basic feeding strategies, or ecological guilds, for feeding at sea: surface feeding, pursuit diving, plunge-diving, and predation of higher vertebrates; within these guilds, there are multiple variations on the theme.
Surface feeding
Many seabirds feed on the ocean's surface, as the action of marine currents often concentrates food such as krill, forage fish, squid, or other prey items within reach of a dipped head.
Surface feeding itself can be broken up into two different approaches, surface feeding while flying (for example as practiced by gadfly petrels, frigatebirds, and storm petrels), and surface feeding while swimming (examples of which are practiced by gulls, fulmars, many of the shearwaters and gadfly petrels). Surface feeders in flight include some of the most acrobatic of seabirds, which either snatch morsels from the water (as do frigate-birds and some terns), or "walk", pattering and hovering on the water's surface, as some of the storm-petrels do. Many of these do not ever land in the water, and some, such as the frigatebirds, have difficulty getting airborne again should they do so. Another seabird family that does not land while feeding is the skimmer, which has a unique fishing method: flying along the surface with the lower mandible in the water—this shuts automatically when the bill touches something in the water. The skimmer's bill reflects its unusual lifestyle, with the lower mandible uniquely being longer than the upper one.
Surface feeders that swim often have unique bills as well, adapted for their specific prey. Prions have special bills with filters called lamellae to filter out plankton from mouthfuls of water, and many albatrosses and petrels have hooked bills to snatch fast-moving prey. On the other hand, most gulls are versatile and opportunistic feeders who will eat a wide variety of prey, both at sea and on land.
Pursuit diving
Pursuit diving exerts greater pressures (both evolutionary and physiological) on seabirds, but the reward is a greater area in which to feed than is available to surface feeders. Underwater propulsion is provided by wings (as used by penguins, auks, diving petrels and some other species of petrel) or feet (as used by cormorants, grebes, loons and several types of fish-eating ducks). Wing-propelled divers are generally faster than foot-propelled divers. The use of wings or feet for diving has limited their utility in other situations: loons and grebes walk with extreme difficulty (if at all), penguins cannot fly, and auks have sacrificed flight efficiency in favour of diving. For example, the razorbill (an Atlantic auk) requires 64% more energy to fly than a petrel of equivalent size. Many shearwaters are intermediate between the two, having longer wings than typical wing-propelled divers but heavier wing loadings than the other surface-feeding procellariids, leaving them capable of diving to considerable depths while still being efficient long-distance travellers. The short-tailed shearwater is the deepest diver of the shearwaters, having been recorded diving below .
Some albatross species are also capable of limited diving, with light-mantled sooty albatrosses holding the record at . Of all the wing-propelled pursuit divers, the most efficient in the air are the albatrosses, and they are also the poorest divers. This is the dominant guild in polar and subpolar environments, but it is energetically inefficient in warmer waters. With their poor flying ability, many wing-propelled pursuit divers are more limited in their foraging range than other guilds.
Plunge diving
Gannets, boobies, tropicbirds, some terns, and brown pelicans all engage in plunge diving, taking fast-moving prey by diving into the water from the flight. Plunge diving allows birds to use the energy from the momentum of the dive to combat natural buoyancy (caused by air trapped in plumage), and thus uses less energy than the dedicated pursuit divers, allowing them to utilise more widely distributed food resources, for example, in impoverished tropical seas. In general, this is the most specialised method of hunting employed by seabirds; other non-specialists (such as gulls and skuas) may employ it but do so with less skill and from lower heights. In brown pelicans, the skills of plunge-diving take several years to fully develop—once mature, they can dive from above the water's surface, shifting the body before impact to avoid injury.
It may be that plunge divers are restricted in their hunting grounds to clear waters that afford a view of their prey from the air. While they are the dominant guild in the tropics, the link between plunge diving and water clarity is inconclusive. Some plunge divers (as well as some surface feeders) are dependent on dolphins and tuna to push shoaling fish up towards the surface.
Kleptoparasitism, scavenging and predation
This catch-all category refers to other seabird strategies that involve the next trophic level up. Kleptoparasites are seabirds that make a part of their living stealing food of other seabirds. Most famously, frigatebirds and skuas engage in this behaviour, although gulls, terns and other species will steal food opportunistically. The nocturnal nesting behaviour of some seabirds has been interpreted as arising due to pressure from this aerial piracy. Kleptoparasitism is not thought to play a significant part of the diet of any species, and is instead a supplement to food obtained by hunting. A study of great frigatebirds stealing from masked boobies estimated that the frigatebirds could at most obtain 40% of the food they needed, and on average obtained only 5%. Many species of gull will feed on seabird and sea mammal carrion when the opportunity arises, as will giant petrels. Some species of albatross also engage in scavenging: an analysis of regurgitated squid beaks has shown that many of the squid eaten are too large to have been caught alive, and include mid-water species likely to be beyond the reach of albatrosses. Some species will also feed on other seabirds; for example, gulls, skuas and pelicans will often take eggs, chicks and even small adult seabirds from nesting colonies, while the giant petrels can kill prey up to the size of small penguins and seal pups.
Life history
Seabirds' life histories are dramatically different from those of land birds. In general, they are K-selected, live much longer (anywhere between twenty and sixty years), delay breeding for longer (for up to ten years), and invest more effort into fewer young. Most species will only have one clutch a year, unless they lose the first (with a few exceptions, like the Cassin's auklet), and many species (like the tubenoses and sulids) will only lay one egg a year.
Care of young is protracted, extending for as long as six months, among the longest for birds. For example, once common guillemot chicks fledge, they remain with the male parent for several months at sea. The frigatebirds have the longest period of parental care of any bird except a few raptors and the southern ground hornbill, with each chick fledging after four to six months and continued assistance after that for up to fourteen months. Due to the extended period of care, breeding occurs every two years rather than annually for some species. This life-history strategy has probably evolved both in response to the challenges of living at sea (collecting widely scattered prey items), the frequency of breeding failures due to unfavourable marine conditions, and the relative lack of predation compared to that of land-living birds.
Because of the greater investment in raising the young and because foraging for food may occur far from the nest site, in all seabird species except the phalaropes, both parents participate in caring for the young, and pairs are typically at least seasonally monogamous. Many species, such as gulls, auks and penguins, retain the same mate for several seasons, and many petrel species mate for life. Albatrosses and procellariids, which mate for life, take many years to form a pair bond before they breed, and the albatrosses have an elaborate breeding dance that is part of pair-bond formation.
Breeding and colonies
Ninety-five percent of seabirds are colonial, and seabird colonies are among the largest bird colonies in the world, providing one of Earth's great wildlife spectacles. Colonies of over a million birds have been recorded, both in the tropics (such as Kiritimati in the Pacific) and in the polar latitudes (as in Antarctica). Seabird colonies occur exclusively for the purpose of breeding; non-breeding birds will only collect together outside the breeding season in areas where prey species are densely aggregated.
Seabird colonies are highly variable. Individual nesting sites can be widely spaced, as in an albatross colony, or densely packed as with a murre colony. In most seabird colonies, several different species will nest on the same colony, often exhibiting some niche separation. Seabirds can nest in trees (if any are available), on the ground (with or without nests), on cliffs, in burrows under the ground and in rocky crevices. Competition can be strong both within species and between species, with aggressive species such as sooty terns pushing less dominant species out of the most desirable nesting spaces. The tropical Bonin petrel nests during the winter to avoid competition with the more aggressive wedge-tailed shearwater. When the seasons overlap, the wedge-tailed shearwaters will kill young Bonin petrels in order to use their burrows.
Many seabirds show remarkable site fidelity, returning to the same burrow, nest or site for many years, and they will defend that site from rivals with great vigour. This increases breeding success, provides a place for returning mates to reunite, and reduces the costs of prospecting for a new site. Young adults breeding for the first time usually return to their natal colony, and often nest close to where they hatched. This tendency, known as philopatry, is so strong that a study of Laysan albatrosses found that the average distance between hatching site and the site where a bird established its own territory was ; another study, this time on Cory's shearwaters nesting near Corsica, found that of nine out of 61 male chicks that returned to breed at their natal colony bred in the burrow they were raised in, and two actually bred with their own mother.
Colonies are usually situated on islands, cliffs or headlands, which land mammals have difficulty accessing. This is thought to provide protection to seabirds, which are often very clumsy on land. Coloniality often arises in types of bird that do not defend feeding territories (such as swifts, which have a very variable prey source); this may be a reason why it arises more frequently in seabirds. There are other possible advantages: colonies may act as information centres, where seabirds returning to the sea to forage can find out where prey is by studying returning individuals of the same species. There are disadvantages to colonial life, particularly the spread of disease. Colonies also attract the attention of predators, principally other birds, and many species attend their colonies nocturnally to avoid predation. Birds from different colonies often forage in different areas to avoid competition.
Migration
Like many birds, seabirds often migrate after the breeding season. Of these, the trip taken by the Arctic tern is the farthest of any bird, crossing the equator in order to spend the Austral summer in Antarctica. Other species also undertake trans-equatorial trips, both from the north to the south, and from south to north. The population of elegant terns, which nest off Baja California, splits after the breeding season with some birds travelling north to the Central Coast of California and some travelling as far south as Peru and Chile to feed in the Humboldt Current. The sooty shearwater undertakes an annual migration cycle that rivals that of the Arctic tern; birds that nest in New Zealand and Chile and spend the northern summer feeding in the North Pacific off Japan, Alaska and California, an annual round trip of .
Other species also migrate shorter distances away from the breeding sites, their distribution at sea determined by the availability of food. If oceanic conditions are unsuitable, seabirds will emigrate to more productive areas, sometimes permanently if the bird is young. After fledging, juvenile birds often disperse further than adults, and to different areas, so are commonly sighted far from a species' normal range. Some species, such as the auks, do not have a concerted migration effort, but drift southwards as the winter approaches. Other species, such as some of the storm petrels, diving petrels and cormorants, never disperse at all, staying near their breeding colonies year round.
Away from the sea
While the definition of seabirds suggests that the birds in question spend their lives on the ocean, many seabird families have many species that spend some or even most of their lives inland away from the sea. Most strikingly, many species breed tens, hundreds or even thousands of miles inland. Some of these species still return to the ocean to feed; for example, the snow petrel, the nests of which have been found inland on the Antarctic mainland, are unlikely to find anything to eat around their breeding sites. The marbled murrelet nests inland in old growth forest, seeking huge conifers with large branches to nest on. Other species, such as the California gull, nest and feed inland on lakes, and then move to the coasts in the winter. Some cormorant, pelican, gull and tern species have individuals that never visit the sea at all, spending their lives on lakes, rivers, swamps and, in the case of some of the gulls, cities and agricultural land. In these cases, it is thought that these terrestrial or freshwater birds evolved from marine ancestors. Some seabirds, principally those that nest in tundra, as skuas and phalaropes do, will migrate over land as well.
The more marine species, such as petrels, auks and gannets, are more restricted in their habits, but are occasionally seen inland as vagrants. This most commonly happens to young inexperienced birds, but can happen in great numbers to exhausted adults after large storms, an event known as a wreck.
Relationship with humans
Seabirds and fisheries
Seabirds have had a long association with both fisheries and sailors, and both have drawn benefits and disadvantages from the relationship.
Fishermen have traditionally used seabirds as indicators of both fish shoals, underwater banks that might indicate fish stocks, and of potential landfall. In fact, the known association of seabirds with land was instrumental in allowing the Polynesians to locate tiny landmasses in the Pacific. Seabirds have provided food for fishermen away from home, as well as bait. Famously, tethered cormorants have been used to catch fish directly. Indirectly, fisheries have also benefited from guano from colonies of seabirds acting as fertilizer for the surrounding seas.
Negative effects on fisheries are mostly restricted to raiding by birds on aquaculture, although long-lining fisheries also have to deal with bait stealing. There have been claims of prey depletion by seabirds of fishery stocks, and while there is some evidence of this, the effects of seabirds are considered smaller than that of marine mammals and predatory fish (like tuna).
Some seabird species have benefited from fisheries, particularly from discarded fish and offal. These discards compose 30% of the food of seabirds in the North Sea, for example, and compose up to 70% of the total food of some seabird populations. This can have other impacts; for example, the spread of the northern fulmar through the United Kingdom is attributed in part to the availability of discards. Discards generally benefit surface feeders, such as gannets and petrels, to the detriment of pursuit divers like penguins and guillemots, which can get entangled in the nets.
Fisheries also have negative effects on seabirds, and these effects, particularly on the long-lived and slow-breeding albatrosses, are a source of increasing concern to conservationists. The bycatch of seabirds entangled in nets or hooked on fishing lines has had a big impact on seabird numbers; for example, an estimated 100,000 albatrosses are hooked and drown each year on tuna lines set out by long-line fisheries. Overall, many hundreds of thousands of birds are trapped and killed each year, a source of concern for some of the rarest species (for example, only about 2,000 short-tailed albatrosses are known to still exist). Seabirds are also thought to suffer when overfishing occurs. Changes to the marine ecosystems caused by dredging, which alters the biodiversity of the seafloor, can also have a negative impact.
Exploitation
The hunting of seabirds and the collecting of seabird eggs have contributed to the declines of many species, and the extinction of several, including the great auk and the spectacled cormorant. Seabirds have been hunted for food by coastal peoples throughout history—one of the earliest instances known is in southern Chile, where archaeological excavations in middens has shown hunting of albatrosses, cormorants and shearwaters from 5000 BP. This pressure has led to some species becoming extinct in many places; in particular, at least 20 species of an original 29 no longer breed on Easter Island. In the 19th century, the hunting of seabirds for fat deposits and feathers for the millinery trade reached industrial levels. Muttonbirding (harvesting shearwater chicks) developed as important industries in both New Zealand and Tasmania, and the name of one species, the providence petrel, is derived from its seemingly miraculous arrival on Norfolk Island where it provided a windfall for starving European settlers. In the Falkland Islands, hundreds of thousands of penguins were harvested for their oil each year. Seabird eggs have also long been an important source of food for sailors undertaking long sea voyages, as well as being taken when settlements grow in areas near a colony. Eggers from San Francisco took almost half a million eggs a year from the Farallon Islands in the mid-19th century, a period in the islands' history from which the seabird species are still recovering.
Both hunting and egging continue today, although not at the levels that occurred in the past, and generally in a more controlled manner. For example, the Māori of Stewart Island / Rakiura continue to harvest the chicks of the sooty shearwater as they have done for centuries, using traditional stewardship, kaitiakitanga, to manage the harvest, but now also work with the University of Otago in studying the populations. In Greenland, however, uncontrolled hunting is pushing many species into steep decline.
Other threats
Other human factors have led to declines and even extinctions in seabird populations and species. Of these, perhaps the most serious are introduced species. Seabirds, breeding predominantly on small isolated islands, are vulnerable to predators because they have lost many behaviours associated with defence from predators. Feral cats can take seabirds as large as albatrosses, and many introduced rodents, such as the Pacific rat, take eggs hidden in burrows. Introduced goats, cattle, rabbits and other herbivores can create problems, particularly when species need vegetation to protect or shade their young. The disturbance of breeding colonies by humans is often a problem as well—visitors, even well-meaning tourists, can flush brooding adults off a colony, leaving chicks and eggs vulnerable to predators.
The build-up of toxins and pollutants in seabirds is also a concern. Seabirds, being apex predators, suffered from the ravages of the insecticide DDT until it was banned; DDT was implicated, for example, in embryo development problems and the skewed sex ratio of western gulls in southern California. Oil spills are also a threat to seabirds: the oil is toxic, and bird feathers become saturated by the oil, causing them to lose their waterproofing. Oil pollution in particular threatens species with restricted ranges or already depressed populations.
Climate change mainly affect seabirds via changes to their habitat: various processes in the ocean lead to decreased availability of food and colonies are more often flooded as a consequence of sea level rise and extreme rainfall events. Heat stress from extreme temperatures is an additional threat. Some seabirds have used changing wind patterns to forage further and more efficiently.
In 2023, plasticosis, a new disease caused solely by plastics, was discovered in seabirds.
The birds identified as having the disease have scarred digestive tracts from ingesting plastic waste. "When birds ingest small pieces of plastic, they found, it inflames the digestive tract. Over time, the persistent inflammation causes tissues to become scarred and disfigured, affecting digestion, growth and survival."
Conservation
The threats faced by seabirds have not gone unnoticed by scientists or the conservation movement. As early as 1903, U.S. President Theodore Roosevelt was convinced of the need to declare Pelican Island in Florida a National Wildlife Refuge to protect the bird colonies (including the nesting brown pelicans), and in 1909 he protected the Farallon Islands. Today many important seabird colonies are given some measure of protection, from Heron Island in Australia to Triangle Island in British Columbia.
Island restoration techniques, pioneered by New Zealand, enable the removal of exotic invaders from increasingly large islands. Feral cats have been removed from Ascension Island, Arctic foxes from many islands in the Aleutian Islands, and rats from Campbell Island. The removal of these introduced species has led to increases in numbers of species under pressure and even the return of extirpated ones. After the removal of cats from Ascension Island, seabirds began to nest there again for the first time in over a hundred years.
Seabird mortality caused by long-line fisheries can be greatly reduced by techniques such as setting long-line bait at night, dying the bait blue, setting the bait underwater, increasing the amount of weight on lines and by using bird scarers, and their deployment is increasingly required by many national fishing fleets.
One of the Millennium Projects in the UK was the Scottish Seabird Centre, near the important bird sanctuaries on Bass Rock, Fidra and the surrounding islands. The area is home to huge colonies of gannets, puffins, skuas and other seabirds. The centre allows visitors to watch live video from the islands as well as learn about the threats the birds face and how we can protect them, and has helped to significantly raise the profile of seabird conservation in the UK. Seabird tourism can provide income for coastal communities as well as raise the profile of seabird conservation, although it needs to be managed to ensure it does not harm the colonies and nesting birds. For example, the northern royal albatross colony at Taiaroa Head in New Zealand attracts 40,000 visitors a year.
The plight of albatross and large seabirds, as well as other marine creatures, being taken as bycatch by long-line fisheries, has been addressed by a large number of non-governmental organizations (including BirdLife International, the American Bird Conservancy and the Royal Society for the Protection of Birds). This led to the Agreement on the Conservation of Albatrosses and Petrels, a legally binding treaty designed to protect these threatened species, which has been ratified by thirteen countries as of 2021 (Argentina, Australia, Brazil, Chile, Ecuador, France, New Zealand, Norway, Peru, South Africa, Spain, Uruguay, United Kingdom).
Role in culture
Many seabirds are little studied and poorly known because they live far out at sea and breed in isolated colonies. Some seabirds, particularly the albatrosses and gulls, are more well known to humans. The albatross has been described as "the most legendary of birds", and have a variety of myths and legends associated with them. While it is widely considered unlucky to harm them, the notion that sailors believed that is a myth that derives from Samuel Taylor Coleridge's famous poem, "The Rime of the Ancient Mariner", in which a sailor is punished for killing an albatross by having to wear its corpse around his neck. Sailors did, however, consider it unlucky to touch a storm petrel, especially one that landed on the ship.
Gulls are one of the most commonly seen seabirds because they frequent human-made habitats (such as cities and dumps) and often show a fearless nature. Gulls have been used as metaphors, as in Jonathan Livingston Seagull by Richard Bach, or to denote a closeness to the sea; in The Lord of the Rings, they appear in the insignia of Gondor and therefore Númenor (used in the design of the films), and they call Legolas to (and across) the sea. Pelicans have long been associated with mercy and altruism because of an early Christian myth that they split open their breast to feed their starving chicks.
Seabird families
The following are the groups of birds normally classed as seabirds. For each order, the species counts given are for only the seabird portions (i.e. the listed groups), not the total number of species.
Sphenisciformes (18 species; Antarctic and southern waters)
Spheniscidae: penguins
Procellariiformes (149 species; pan-oceanic and pelagic)
Diomedeidae: albatrosses
Procellariidae: petrels (including fulmars, prions, shearwaters, gadfly petrels, diving petrels, and other petrels)
Hydrobatidae: northern storm petrels
Oceanitidae: southern storm petrels
Pelecaniformes (8 species; worldwide)
Pelecanidae: pelicans
Suliformes (57 species; worldwide)
Sulidae: gannets and boobies
Phalacrocoracidae: cormorants
Fregatidae: frigatebirds
Phaethontiformes (3 species; worldwide tropical seas)
Phaethontidae: tropicbirds
Charadriiformes (138 species; worldwide)
Laridae: larids (including gulls, terns, and skimmers)
Stercorariidae: skuas
Alcidae: auks
Genus Phalaropus within Scolopacidae: phalaropes
For an alternative taxonomy of these groups, see also Sibley-Ahlquist taxonomy.
| Biology and health sciences | General articles | null |
147020 | https://en.wikipedia.org/wiki/Hygiene | Hygiene | Hygiene is a set of practices performed to preserve health.
According to the World Health Organization (WHO), "Hygiene refers to conditions and practices that help to maintain health and prevent the spread of diseases." Personal hygiene refers to maintaining the body's cleanliness. Hygiene activities can be grouped into the following: home and everyday hygiene, personal hygiene, medical hygiene, sleep hygiene, and food hygiene. Home and every day hygiene includes hand washing, respiratory hygiene, food hygiene at home, hygiene in the kitchen, hygiene in the bathroom, laundry hygiene, and medical hygiene at home. And also environmental hygiene in the society to prevent all kinds of bacterias from penetrating into our homes.
Many people equate hygiene with "cleanliness", but hygiene is a broad term. It includes such personal habit choices as how frequently to take a shower or bath, wash hands, trim fingernails, and wash clothes. It also includes attention to keeping surfaces in the home and workplace clean, including bathroom facilities. Adherence to regular hygiene practices is often regarded as a socially responsible and respectable behavior, while neglecting proper hygiene can be perceived as unclean or unsanitary, and may be considered socially unacceptable or disrespectful, while also posing a risk to public health.
Definition and overview
Hygiene is a practice related to lifestyle, cleanliness, health, and medicine. In medicine and everyday life, hygiene practices are preventive measures that reduce the incidence and spread of germs leading to disease.
Hygiene practices vary from one culture to another.
In the manufacturing of food, pharmaceuticals, cosmetics, and other products, good hygiene is a critical component of quality assurance.
The terms cleanliness and hygiene are often used interchangeably, which can cause confusion. In general, hygiene refers to practices that prevent spread of disease-causing organisms. Cleaning processes (e.g., handwashing) remove infectious microbes as well as dirt and soil, and are thus often the means to achieve hygiene.
Other uses of the term are as follows: body hygiene, personal hygiene, sleep hygiene, mental hygiene, dental hygiene, and occupational hygiene, used in connection with public health.
Home hygiene overview
Home hygiene pertains to the hygiene practices that prevent or minimize the spread of disease at home and other everyday settings such as social settings, public transport, the workplace, public places, and more. Hygiene in a variety of settings plays an important role in preventing the spread of infectious diseases. It includes procedures like hand hygiene, respiratory hygiene, food and water hygiene, general home hygiene (hygiene of environmental sites and surfaces), care of domestic animals, and home health care (the care of those who are at greater risk of infection).
At present, these components of hygiene tend to be regarded as separate issues, although based on the same underlying microbiological principles. Preventing the spread of diseases means breaking the chain of infection transmission so that infection cannot spread. "Targeted hygiene" is based on identifying the routes of pathogen spread in the home and introducing hygiene practices at critical times to break the chain of infection. It uses a risk-based approach based on Hazard Analysis Critical Control Point (HACCP).
The main sources of infection in the home are people (who are carriers or are infected), foods (particularly raw foods), water, pets, and domestic animals. Sites that accumulate stagnant water – such as sinks, toilets, waste pipes, cleaning tools, and face cloths – readily support microbial growth and can become secondary reservoirs of infection, though species are mostly those that threaten "at risk" groups. Pathogens (such as potentially infectious bacteria and viruses – colloquially called "germs") are constantly shed via mucous membranes, feces, vomit, skin scales, and other means. When circumstances combine, people are exposed, either directly or via food or water, and can develop an infection.
The main "highways" for the spread of pathogens in the home are the hands, hand and food contact surfaces, and cleaning cloths and utensils (e.g. fecal–oral route of transmission). Pathogens can also be spread via clothing and household linens, such as towels. Utilities such as toilets and wash basins were invented to deal safely with human waste but still have risks associated with them. Safe disposal of human waste is a fundamental need; poor sanitation is a primary cause of diarrhea disease in low-income communities. Respiratory viruses and fungal spores spread via the air.
Good home hygiene means engaging in hygiene practices at critical points to break the chain of infection. Because the "infectious dose" for some pathogens can be very small (10–100 viable units or even less for some viruses), and infection can result from direct transfer of pathogens from surfaces via hands or food to the mouth, nasal mucous, or the eye, "hygienic cleaning" procedures should be adopted to eliminate pathogens from critical surfaces.
Hand washing
Respiratory hygiene
Correct respiratory and hand hygiene when coughing and sneezing reduces the spread of pathogens particularly during the cold and flu season:
Carry tissues and use them to catch coughs and sneezes, or sneeze into your elbow.
Dispose of tissues as soon as possible.
Hygiene in the kitchen, bathroom and toilet
Routine cleaning of hands, food, sites, and surfaces (such as toilet seats and flush handles, door and tap handles, work surfaces, and bath and basin surfaces) in the kitchen, bathroom, and toilet rooms reduces the spread of pathogens. The infection risk from flush toilets is not high, provided they are properly maintained, although some splashing and aerosol formation can occur during flushing, particularly when someone has diarrhea. Pathogens can survive in the scum or scale left behind on baths, showers, and washbasins after washing and bathing.
Thorough cleaning is important to prevent the spread of fungal infections. Molds can live on wall and floor tiles and on shower curtains. Mold can be responsible for infections, cause allergic reactions, deteriorate/damage surfaces, and cause unpleasant odors. Primary sites of fungal growth are inanimate surfaces, including carpets and soft furnishings. Airborne fungi are usually associated with damp conditions, poor ventilation, or closed air systems.
Hygienic cleaning can be done through:
Mechanical removal (i.e., cleaning) using a soap or detergent. To be effective as a hygiene measure, this process must be followed by thorough rinsing under running water to remove pathogens from the surface.
Using a process or product that inactivates the pathogens in situ. Pathogen kill is achieved using a "micro-biocidal" product, i.e., a disinfectant or antibacterial product; waterless hand sanitizer; or by application of heat.
In some cases, combined pathogen removal with kill is used, e.g., laundering of clothing and household linens such as towels and bed linen.
House deep-cleaning an intensive cleaning process targeting often-neglected areas, enhancing aesthetics, and improving health by reducing allergens and bacteria. It typically includes tasks like detailed dusting, appliance cleaning, and carpet shampooing, recommended biannually to maintain a home's hygiene and air quality.
Laundry hygiene
Laundry hygiene involves practices that prevent disease and its spread via soiled clothing and household linens such as towels. Items most likely to be contaminated with pathogens are those that come into direct contact with the body, e.g., underwear, personal towels, facecloths, nappies. Cloths or other fabric items used during food preparation, or for cleaning the toilet or cleaning up material such as feces or vomit are a particular risk.
Microbiological and epidemiological data indicates that clothing and household linens are a risk factor for infection transmission in home and everyday life settings as well as institutional settings. The lack of quantitative data linking contaminated clothing to infection in the domestic setting makes it difficult to assess the extent of this risk. This also indicates that risks from clothing and household linens are somewhat less than those associated with hands, hand contact and food contact surfaces, and cleaning cloths, but even so these risks need to be managed through effective laundering practices. In the home, this should be carried out as part of a multibarrier approach to hygiene which includes hand, food, respiratory, and other hygiene practices.
Infectious disease risks from contaminated clothing can increase significantly under certain conditions - for example, in healthcare situations in hospitals, care homes, and the domestic setting where someone has diarrhoea, vomiting, or a skin or wound infection. The risk increases in circumstances where someone has reduced immunity to infection.
Hygiene measures, including laundry hygiene, are an important part of reducing spread of antibiotic-resistant strains of infectious organisms. In the community, otherwise-healthy people can become persistent skin carriers of MRSA, or faecal carriers of enterobacteria strains which can carry multi-antibiotic resistance factors (e.g. NDM-1 or ESBL-producing strains). The risks are not apparent until, for example, they are admitted to hospital, when they can become "self infected" with their own resistant organisms following a surgical procedure. As persistent nasal, skin, or bowel carriage in the healthy population spreads "silently" across the world, the risks from resistant strains in both hospitals and the community increases. In particular the data indicates that clothing and household linens are a risk factor for spread of S. aureus (including MRSA and PVL-producing MRSA strains), and that effectiveness of laundry processes may be an important factor in defining the rate of community spread of these strains. Experience in the United States suggests that these strains are transmissible within families and in community settings such as prisons, schools, and sport teams. Skin-to-skin contact (including unabraded skin) and indirect contact with contaminated objects such as towels, sheets, and sports equipment seem to represent the mode of transmission.
During laundering, temperature and detergent work to reduce microbial contamination levels on fabrics. Soil and microbes from fabrics are severed and suspended in the wash water. These are then "washed away" during the rinse and spin cycles. In addition to physical removal, micro-organisms can be killed by thermal inactivation which increases as the temperature is increased. Chemical inactivation of microbes by the surfactants and activated oxygen-based bleach used in detergents contributes to the hygiene effectiveness of laundering. Adding hypochlorite bleach in the washing process achieves inactivation of microbes. A number of other factors can contribute including drying and ironing.
Drying laundry on a line in direct sunlight is known to reduce pathogens.
In 2013, the International Scientific Forum on Home Hygiene reviewed 30 studies of the hygiene effectiveness of laundering at temperatures ranging from room temperature to , under varying conditions. A key finding was the lack of standardization and control within studies, and the variability in test conditions between studies such as wash cycle time, number of rinses, and other factors. The consequent variability in the data (i.e., the reduction in contamination on fabrics) in turn makes it extremely difficult to propose guidelines for laundering with any confidence. As a result, there is significant variability in the recommendations for hygienic laundering given by different agencies.
Medical hygiene at home
Medical hygiene pertains to hygiene practices that prevent or minimize disease and the spreading of disease in relation to administering medical care to those who are infected or who are more at risk of infection in the home. Members of "at-risk" groups are cared for at home by a carer who may be a household member and who requires a good knowledge of hygiene. People with reduced immunity to infection, who are looked after at home, make up an increasing proportion of the population (, up to 20%). The largest proportion are the elderly who have co-morbidities that reduce their immunity to infection. It also includes the very young, patients discharged from hospital, taking immuno-suppressive drugs, or using invasive systems, etc. For patients discharged from hospital, or being treated at home, special "medical hygiene" procedures may need to be performed for them, such as catheter or dressing replacement, which puts them at higher risk of infection.
Antiseptics may be applied to cuts, wounds, and abrasions of the skin to prevent the entry of harmful bacteria that can cause sepsis. Day-to-day hygiene practices, other than special medical hygiene procedures, are no different for those at increased risk of infection than for other family members. The difference is that, if hygiene practices are not correctly carried out, the risk of infection is much greater.
Disinfectants and antibacterials in home hygiene
Chemical disinfectants are products that kill pathogens. If the product is a disinfectant, the label on the product should say "disinfectant" or "kills" pathogens. Some commercial products, e.g. bleaches, even though they are technically disinfectants, say that they "kill pathogens" but are not actually labelled as "disinfectants". Not all disinfectants kill all types of pathogens. All disinfectants kill bacteria (called bactericidal). Some also kill fungi (fungicidal), bacterial spores (sporicidal), or viruses (virucidal).
An antibacterial product acts against bacteria in some unspecified way. Some products labelled "antibacterial" kill bacteria while others may contain a concentration of active ingredient that only prevents them from multiplying. It is, therefore, important to check whether the product label states that it "kills bacteria". An antibacterial is not necessarily anti-fungal or anti-viral unless this is stated on the label.
The term sanitizer has been used to define substances that both clean and disinfect. More recently this term has been applied to alcohol-based products that disinfect the hands (alcohol hand sanitizers). Alcohol hand sanitizers however are not considered to be effective on soiled hands.
The term biocide is a broad term for a substance that kills, inactivates or otherwise controls living organisms. It includes antiseptics and disinfectants, which combat micro-organisms, and pesticides.
Personal hygiene
Regular activities
Personal hygiene involves those practices performed by a person to care for their bodily health and well-being through cleanliness. Motivations for personal hygiene practice include reduction of personal illness, healing from illness, optimal health and sense of wellbeing, social acceptance, and prevention of spread of illness to others. What is considered proper personal hygiene can be culture-specific and may change over time.
Practices that are generally considered proper hygiene include showering or bathing regularly, washing hands regularly and especially before handling food, face washing, washing scalp hair, keeping hair short or removing hair, wearing clean clothing, brushing teeth, and trimming fingernails and toenails. Some practices are sex-specific, such as by a woman during menstruation.
Toiletry bags hold body hygiene and toiletry supplies.
Anal hygiene is the practice that a person performs on their anal area after defecation. The anus and buttocks may be either washed with liquids or wiped with toilet paper, or by adding gel wipe to toilet tissue as an alternative to wet wipes or other solid materials in order to remove remnants of feces.
People tend to develop a routine for attending to their personal hygiene needs. Other personal hygienic practices include covering one's mouth when coughing, disposal of soiled tissues appropriately, making sure toilets are clean, and making sure food handling areas are clean, besides other practices. Some cultures do not kiss or shake hands in order to reduce transmission of bacteria by contact.
Personal grooming extends personal hygiene as it pertains to the maintenance of a good personal and public appearance, which need not necessarily be hygienic. It may involve, for example, using deodorants or perfume, shaving, or combing.
Hygiene of internal ear canals
Excessive cleaning of the ear canals can result in infection or irritation. The ear canals require less care than other parts of the body because they are sensitive and mostly self-cleaning. There is a slow and orderly migration of the skin lining the ear canal from the eardrum to the outer opening of the ear. Old earwax is constantly being transported from the deeper areas of the ear canal out to the opening where it usually dries, flakes, and falls out. Attempts to clean the ear canals through the removal of earwax can push debris and foreign material into the ear that the natural movement of ear wax out of the ear would have removed.
Oral hygiene
It is recommended that all healthy adults brush twice a day, softly, with the correct technique, replacing their toothbrush every few months (~3).
There are a number of common oral hygiene misconceptions. The National Health Service (NHS) of England recommends not rinsing the mouth with water after brushing – only to spit out excess toothpaste. They claim that this helps fluoride from toothpaste bond to teeth for its preventative effects against tooth decay. It is also not recommended to brush immediately after drinking acidic substances, including sparkling water. It is also recommended to floss once a day, with a different piece of floss at each flossing session. The effectiveness of amorphous calcium phosphate products, such as Tooth Mousse, is in debate. Visits to a dentist for a checkup every year at least are recommended.
Sleep hygiene
Sleep hygiene is the recommended behavioral and environmental practices that promote better quality sleep. These recommendations were developed in the late 1970s as a method to help people with mild to moderate insomnia, but, , the evidence for effectiveness of individual recommendations is "limited and inconclusive". Clinicians assess the sleep hygiene of people who present with insomnia and other conditions, such as depression, and offer recommendations based on the assessment. Sleep hygiene recommendations include establishing a regular sleep schedule, using naps with care, not exercising physically or mentally too close to bedtime, and avoiding alcohol as well as nicotine, caffeine, and other stimulants in the hours before bedtime. Further recommendations include limiting worry, limiting exposure to light in the hours before sleep, getting out of bed if sleep does not come, not using the bed for anything but sleep, and having a peaceful, comfortable, and dark sleep environment.
Personal care services hygiene
Personal care services hygiene pertains to the care and use of instruments used in the administration of personal care services to people:
Personal care hygiene practices include:
sterilization of instruments used by service providers including hairdressers, aestheticians, and other service providers
sterilization by autoclave of instruments used in body piercing and tattooing
cleaning hands
Challenges
Excessive body hygiene is a possible sign of obsessive–compulsive disorder. Neglecting bodily hygiene, or the cleanliness of one's environment, may be a sign of major depression and other psychological disorders.
Hygiene hypothesis and allergies
Although media coverage of the hygiene hypothesis has declined, popular folklore continues to sometimes assert that dirt is healthy and hygiene unnatural. This has caused health professionals to be concerned that hygiene behaviors which are the foundation of public health are being undermined. In response to the need for effective hygiene in home and everyday life settings, the International Scientific Forum on Home Hygiene developed a "risk-based" or targeted approach to home hygiene that seeks to ensure that hygiene measures are focused on the places and times most critical for infection transmission. While targeted hygiene was originally developed as an effective approach to hygiene practice, it also seeks, as far as possible, to sustain "normal" levels of exposure to the microbial flora of our environment to the extent that is important to build a balanced immune system.
Although there is substantial evidence that some microbial exposures in early childhood can in some way protect against allergies, there is no evidence that humans need exposure to harmful microbes (infection) or that it is necessary to develop a clinical infection. Nor is there evidence that hygiene measures such as hand washing, food hygiene, etc., are linked to increased susceptibility to atopic disease. If this is the case, there is no conflict between the goals of preventing infection and minimizing allergies. that the answer lies in more fundamental changes in lifestyles that have led to decreased exposure to certain microbial or other species, such as helminths, that are important for development of immuno-regulatory mechanisms. There is still much uncertainty as to which lifestyle factors are involved.
Medical hygiene
Medical hygiene pertains to hygiene practices related to the administration of medicine and medical care that prevents or minimizes the spread of disease.
Medical hygiene practices include:
isolation of infectious persons or materials to prevent spread of infection
sterilization of instruments used in surgical procedures
proper bandaging and dressing of injuries
safe disposal of medical waste
disinfection of reusables (i.e., linen, pads, uniforms)
scrubbing up, handwashing, especially in an operating room, but in more general health-care settings as well, where diseases can be transmitted
ethanol-based sanitizers
Most of these practices were developed in the 19th century and were well-established by the mid-20th century. Some procedures (such as disposal of medical waste) were refined in response to late-20th century disease outbreaks, notably AIDS and Ebola.
Food hygiene
Culinary hygiene (or food hygiene) pertains to practices of food management and cooking that prevent food contamination, prevent food poisoning, and minimize the transmission of disease to other foods, humans, or animals. Culinary hygiene practices specify safe ways to handle, store, prepare, serve, and eat food.
Hygiene aspects in low- and middle-income countries
In developing countries (or low- and middle-income countries), universal access to water and sanitation, coupled with hygiene promotion, is essential in reducing infectious diseases. This approach has been integrated into the Sustainable Development Goal Number 6 whose second target states: "By 2030, achieve access to adequate and equitable sanitation and hygiene for all and end open defecation, paying special attention to the needs of women and girls and those in vulnerable situations". Due to their close linkages, water, sanitation, hygiene are together abbreviated and funded under the term WASH in development cooperation.
About two million people die every year due to diarrheal diseases; most of them are children less than five years of age. The most affected are people in developing countries who live in extreme conditions of poverty, normally peri-urban dwellers or rural inhabitants. Providing access to sufficient quantities of safe water and facilities for a sanitary disposal of excreta, and introducing sound hygiene behaviors are important in order to reduce the burden of disease.
Research shows that, if widely practiced, hand washing with soap could reduce diarrhea by almost fifty percent and respiratory infections by nearly twenty-five percent Hand washing with soap also reduces the incidence of skin diseases, and eye infections like trachoma and intestinal worms, especially ascariasis and trichuriasis. Other hygiene practices, such as safe disposal of waste, surface hygiene, and care of domestic animals, are important in low income communities to break the chain of infection transmission.
Cleaning of toilets and hand wash facilities is important to prevent odors and make them socially acceptable. Social acceptance is an important part of encouraging people to use toilets and wash their hands, in situations where open defecation is still seen as a possible alternative, e.g. in rural areas of some developing countries.
Household water treatment and safe storage
Household water treatment and safe storage ensure drinking water is safe for consumption. These interventions are part of the approach of self-supply of water for households. Drinking water quality remains a significant problem in developing and in developed countries; even in the European region it is estimated that 120 million people do not have access to safe drinking water. Point-of-use water quality interventions can reduce diarrheal disease in communities where water quality is poor or in emergency situations where there is a breakdown in water supply.
Since water can become contaminated during storage at home (e.g. by contact with contaminated hands or using dirty storage vessels), safe storage of water in the home is important.
Methods for treatment of drinking water at the household level include:
chemical disinfection using chlorine or iodine
boiling
filtration using ceramic filters
solar disinfection — Solar disinfection is an effective method, especially when no chemical disinfectants are available.
UV irradiation — Community or household UV systems may be batch or flow-though. The lamps can be suspended above the water channel or submerged in the water flow.
combined flocculation/disinfection systems — available as sachets of powder that act by coagulating and flocculating sediments in water followed by release of chlorine
multibarrier methods — Some systems use two or more of the above treatments in combination or in succession to optimize efficacy.
portable water purification devices
History
Asia
China
Bathing culture in Chinese literature can be traced back to the Shang dynasty (), when Oracle bone inscriptions describe people washing their hair and body in a bath. The Book of Rites, a work regarding Zhou dynasty () ritual, politics, and culture compiled during the Warring States period, recommends that people take a hot shower every five days, and wash their hair every three days. It was also considered good manners to take a bath provided by the host before a dinner. In the Han dynasty, bathing became a regular activity, and for government officials bathing was required every five days.
Ancient bath facilities have been found in ancient Chinese cities, such as Dongzhouyang archaeological site in Henan Province. Bathrooms were called (), and bathtubs were made of bronze or timber. Bath beans – a powdery soap mixture of ground beans, cloves, eaglewood, flowers, and even powdered jade – were recorded in the Han Dynasty. Bath beans were considered luxury toiletries, while common people simply used powdered beans without spices mixed in. Luxurious bathhouses built around hot springs were recorded in Tang dynasty. While royal bathhouses and bathrooms were common among ancient Chinese nobles and commoners, public bathhouses were a relatively late development. In the Song dynasty (), public bathhouses became popular and people could find them readily. Bathing became an essential part of social life and recreation. Bathhouses often provided massage, nail cutting service, rubdown service, ear cleaning, food, and beverages. Marco Polo, who traveled to China during the Yuan dynasty, noted Chinese bathhouses were using coal to heat the bathhouse, which he had never seen before in Europe. Coal was so plentiful that Chinese people of every social class had bathrooms in their houses, and people took showers every day in the winter for enjoyment.
A typical Ming dynasty bathhouse had slabbed floors and brick domed ceilings. A huge boiler would be installed in the back of the house, connected with the bathing pool through a tunnel. Water could be pumped into the pool by turning wheels attended by the staff.
Japan
The origin of Japanese bathing is , ritual purification with water.
In the Heian period (), houses of prominent families, such as the families of court nobles or samurai, had baths. The bath had lost its religious significance and instead became leisure. became (to bathe in a shallow wooden tub). In the 17th century, the first European visitors to Japan recorded the habit of daily baths in mixed sex groups.
Indian subcontinent
The earliest written account of elaborate codes of hygiene can be found in several Hindu texts, such as the Manusmriti and the Vishnu Purana. Bathing is one of the five (daily duties) in Hinduism, and not performing it leads to sin, according to some scriptures.
Ayurveda is a system of medicine developed in ancient times that is still practiced in India, mostly combined with conventional Western medicine. Contemporary Ayurveda stresses a sattvic diet and good digestion and excretion. Hygiene measures include oil pulling, and tongue scraping. Detoxification also plays an important role.
The Americas
Mesoamerica
Spanish chronicles describe the bathing habits of the peoples of Mesoamerica during and after the conquest.
Bernal Díaz del Castillo describes Moctezuma (the Mexica, or Aztec, emperor at the arrival of Cortés) in his as being "...Very neat and cleanly, bathing every day each afternoon...".
Bathing was not restricted to the elite, but was practiced by all people; the chronicler Tomás López Medel wrote after a journey to Central America that " and the custom of washing oneself is so quotidian [common] amongst the Indians, both of cold and hot lands, as is eating, and this is done in fountains and rivers and other water to which they have access, without anything other than pure water..."
The Mesoamerican bath, known as in Spanish, from the Nahuatl word , a compound of ("steam") and ("house"), consists of a room, often in the form of a small dome, with an exterior firebox known as () that heats a small portion of the room's wall made of volcanic rocks; after this wall has been heated, water is poured on it to produce steam, an action known as . As the steam accumulates in the upper part of the room a person in charge uses a bough to direct the steam to the bathers who are lying on the ground, with which he later gives them a massage, then the bathers scrub themselves with a small flat river stone and finally the person in charge introduces buckets with water along with soap and grass used to rinse. This bath had also ritual importance, and was tied to the goddess Toci; it is also therapeutic when medicinal herbs are used in the water for the . It is still used in Mexico.
Europe
Antiquity
Regular bathing was a hallmark of Roman civilization. Elaborate baths were constructed in urban areas to serve the public, who typically demanded the infrastructure to maintain personal cleanliness. The complexes usually consisted of large, swimming pool-like baths, smaller cold and hot pools, saunas, and spa-like facilities where people could be depilated, oiled, and massaged. Water was constantly changed by an aqueduct-fed flow. Bathing outside of urban centers involved smaller, less elaborate bathing facilities, or simply the use of clean bodies of water. Roman cities also had large sewers, such as Rome's Cloaca Maxima, into which public and private latrines drained. Romans did not have demand-flush toilets but did have some toilets with a continuous flow of water under them. The Romans used scented oils (mostly from Egypt), among other alternatives.
Christianity has always placed a strong emphasis on hygiene. Despite rejecting mixed bathing, early Christian clergy encouraged believers to bath, which contributed to hygiene and good health according to the Church Fathers Clement of Alexandria and Tertullian. The Church built public bathing facilities that were separated by sex near monasteries and pilgrimage sites.
Middle Ages
Contrary to popular belief, bathing and sanitation were not lost in Europe with the collapse of the Roman Empire. Starting in the early Middle Ages, popes situated baths within church basilicas and monasteries. Pope Gregory the Great promoted bathing as a bodily need. The use of water in many Christian countries is partly due to Biblical toilet etiquette which encourages washing after all instances of defecation. Bidet and bidet showers were used in regions where water was considered essential for anal cleansing. Public bathhouses were common in medieval Christendom larger towns and cities such as Constantinople, Paris, Regensburg, Rome and Naples. Great bathhouses were built in Byzantine centers such as Constantinople and Antioch.
In the 11th and 12th centuries, bathing was essential to the Western European upper class: the Cluniac monasteries (popular centers for resorting and retiring) were always equipped with bathhouses. These baths were also used ritually when the monks took full immersion baths at the two Christian festivals of renewal. The rules of the Augustinians and Benedictines contained references to ritual purification, and, inspired by Benedict of Nursia, encouraged the practice of therapeutic bathing. Benedictine monks also played a role in the development and promotion of spas.
On the other hand, bathing also sparked erotic phantasies, played upon by the writers of romances intended for the upper class; in the tale of Melusine the bath was a crucial element of the plot.
Cities regulated public bathing – the 26 public baths of Paris in the late 13th century were strictly overseen by the civil authorities and guild laws banned prostitutes from bathhouse admission.
In 14th century Tuscany, newlywed couples commonly took a bath together and we find an illustration of this custom in a fresco in the town hall of San Gimignano.
As evident in Hans Folz' Bath Booklet (a late 15th century guide on European baths) and various artistic depictions such as Albrecht Dürer's Women's Bath , public bathing continued to be a popular past time in the Renaissance. In Britain, the rise of Protestantism also played a prominent role in the development of spa culture.
Modernity
Until the late 19th century, only the elite in Western cities typically possessed indoor facilities for relieving bodily functions. The poorer majority used communal facilities built above cesspools in backyards and courtyards. This changed after Dr. John Snow discovered that cholera was transmitted by the fecal contamination of water. Though it took decades for his findings to gain wide acceptance, governments and sanitary reformers were eventually convinced of the health benefits of using sewers to keep human waste from contaminating the water. This encouraged the widespread adoption of both the flush toilet and the moral imperative that bathrooms should be indoors and as private as possible.
Modern sanitation was not widely adopted until the 19th and 20th centuries. According to medieval historian Lynn Thorndike, people in Medieval Europe probably bathed more than people did in the 19th century. Some time after Louis Pasteur's experiments proved the germ theory of disease and Joseph Lister and others put this into practice in sanitation, hygienic practices came to be regarded as synonymous with health, as they are in modern times.
The importance of hand washing for human healthparticularly for people in vulnerable circumstances like mothers who had just given birth or wounded soldiers in hospitalswas first recognized in the mid 19th century by two pioneers of hand hygiene: the Hungarian physician Ignaz Semmelweis who worked in Vienna, Austria, and Florence Nightingale, the English "founder of modern nursing". At that time most people still believed that infections were caused by foul odors called miasmas.
Middle East
Islam stresses the importance of cleanliness and personal hygiene. Islamic hygienical jurisprudence, which dates back to the 7th century, has a number of elaborate rules. (ritual purity) involves performing (ablution) for the five daily (prayers), as well as regularly performing (bathing), which led to bathhouses being built across the Islamic world. Islamic toilet hygiene also requires washing with water after using the toilet, for purity and to minimize pathogens.
In the Abbasid Caliphate (8th–13th centuries), its capital city of Baghdad (Iraq) had 65,000 baths, along with a sewer system. Cities and towns of the medieval Islamic world had water supply systems powered by hydraulic technology that supplied drinking water along with much greater quantities of water for ritual washing, mainly in mosques and hammams (baths). Bathing establishments in various cities were rated by Arabic writers in travel guides. Medieval Islamic cities such as Baghdad, Córdoba (Islamic Spain), Fez (Morocco), and Fustat (Egypt) also had sophisticated waste disposal and sewage systems with interconnected networks of sewers. The city of Fustat also had multi-storey tenement buildings (with up to six floors) with flush toilets, which were connected to a water supply system, and flues on each floor carrying waste to underground channels.
A basic form of contagion theory dates back to the Persian medicine in the medieval, where it was proposed by Persian physician Ibn Sina (also known as Avicenna) in The Canon of Medicine (1025), the most authoritative medical textbook of the Middle Ages. He mentioned that people can transmit disease to others by breath, noted contagion with tuberculosis, and discussed the transmission of disease through water and dirt. The concept of invisible contagion was eventually widely accepted by Islamic scholars. In the Ayyubid Sultanate, they referred to them as ("impure substances"). The scholar Ibn al-Haj al-Abdari (), while discussing Islamic diet and hygiene, gave advice and warnings about how contagion can contaminate water, food, and garments, and could spread through the water supply.
In the 9th century, Ziryab invented a type of deodorant. He also promoted morning and evening baths, and emphasized the maintenance of personal hygiene. Ziryab is thought to have invented a type of toothpaste, which he popularized throughout Islamic Iberia. The exact ingredients of this toothpaste are not known, but it was reported to have been both "functional and pleasant to taste."
Sub-Saharan Africa
In West Africa, various ethnic groups such as the Yoruba have used black soap to treat skin diseases.
In Southern Africa, the Zulu people conducted methods of sanitation by using water stored in pottery at Ulundi. The Himba people of Namibia and Angola also utilized mixtures of smoke and otjitze treat skin diseases in regions where water is scarce.
Soap and soap makers
Hard toilet soap with a pleasant smell was invented in the Middle East during the Islamic Golden Age when soap-making became an established industry. Recipes for soap-making are described by Muhammad ibn Zakariya al-Razi (), who also gave a recipe for producing glycerine from olive oil. In the Middle East, soap was produced from the interaction of fatty oils and fats with alkali. In Syria, soap was produced using olive oil together with alkali and lime. Soap was exported from Syria to other parts of the Muslim world and to Europe. Two key Islamic innovations in soapmaking was the invention of bar soap, described by al-Razi, and the addition of scents using perfume technology perfected in the Islamic world.
By the 15th century, the manufacture of soap in Christendom had become virtually industrialized, with sources in Antwerp, Castile, Marseille, Naples, and Venice. In the 17th century the Spanish Catholic manufacturers purchased the monopoly on Castile soap from the cash-strapped Carolinian government. Industrially-manufactured bar soaps became available in the late 18th century, as advertising campaigns in Europe and America promoted popular awareness of the relationship between cleanliness and health.
A major contribution of the Christian missionaries in Africa, China, Guatemala, India, Indonesia, Korea, and other places was better health care through hygiene and introducing and distributing soap, and "cleanliness and hygiene became an important marker of being identified as a Christian".
Society and culture
Religious hygienic customs
Many religions require or encourage ritual purification via bathing or immersing the hands in water. In Islam, washing oneself via or is necessary for performing prayer. Islamic tradition also lists a variety of rules concerning proper hygiene after using the bathroom. The Baháʼí Faith mandates the washing of the hands and face prior to the obligatory Baháʼí prayers. Orthodox Judaism requires a bath following menstruation and childbirth, while washing the hands is performed upon waking up and before eating bread. Water plays a role in Christian rituals as well, and in certain denominations of Christianity such as the Ethiopian Orthodox Tewahedo Church which prescribes several kinds of hand washing for example after leaving the latrine, lavatory, or bathhouse, or before prayer, or after eating a meal, or ritual handwashing.
Etymology
First attested in English in 1676, the word hygiene comes from the French , the latinisation of the Greek () , meaning "(art) of health", from , "good for the health, healthy", in turn from (), "healthful, sound, salutary, wholesome". In ancient Greek religion, Hygeia () was the personification of health, cleanliness, and hygiene.
| Biology and health sciences | Health, fitness, and medicine | null |
147190 | https://en.wikipedia.org/wiki/Ground%20stone | Ground stone | In archaeology, ground stone is a category of stone tool formed by the grinding of a coarse-grained tool stone, either purposely or incidentally. Ground stone tools are usually made of basalt, rhyolite, granite, or other cryptocrystalline and igneous stones whose coarse structure makes them ideal for grinding other materials, including plants and other stones.
Organic and inorganic materials are processed on ground stones into edible products. They are sometimes the only artefacts preserved on archaeological sites and are found worldwide.
Origin
The adoption of ground stone technology is associated closely with the Neolithic, also called the New Stone Age. The Stone Age comes from the three-age system developed by Christian Jürgensen Thomsen. In the Levant ground stones appear in Mesolithic 2 (Natufian). In prehistoric Japan, ground stone tools appear during the Japanese Paleolithic, possibly predating adoption elsewhere in the Neolithic by 25,000 years.
Creation
Ground stones were created and used for a wide variety of reasons. Each use resulted in a different development and process by which a person created their ground stone. For example, the process for creating the head of a hammer is different from the process used to create a detailed decoration piece for one’s home. That being said, some processes are basic to most ground stone making.
When choosing what type of stone to use for a ground stone tool, toughness is the most important factor. If the stone is not tough enough to withstand hard hits and instead just flakes and cracks easily, the work done to create the tool has gone to waste. A stone that will not shear, flake, or crack when tested against large impacts is the most important aspect when choosing what kind of stone to use. Examples of this kind of stone include limestone, sandstone, granite, basalt, rhyolite and other igneous and cryptocrystalline rocks.
Cryptocrystalline rocks are good to use for ground stones because they have a very fine grain structure. This is helpful because the smaller the grains are in a rock, the harder the rock is.
Holes could be ground out of stones with the use of sharp pointed stones or hardened sticks. By spinning the ground stone with one's hands and applying substantial pressure to the sharp point into the ground stone, a hole could be drilled into the stone with a large amount of time and effort. Sand would be used to help quicken the process by putting it in the partially formed hole as the sharp point was being pressed. The sand would help grind more of the stone away. To put a hole all the way through a piece of stone, it would be first drilled half way in one direction and be finished on the opposite side.
Some ground stone tools are incidental, caused by use with other tools: manos, for example, are hand stones used in conjunction with metates and other grinding slabs (querns), and develop their ground surfaces through wear. Other ground stone tools include adzes, celts, and axes, which are manufactured using a labor-intensive, time-consuming method of repeated grinding against a harder stone or with sand, often using water as a lubricant. These tools are often made using durable finer-grained materials rather than coarse materials. In the North American arctic, tools made of ground slate were used by the Norton, Dorset, and Thule tool cultures, among others. Common forms of these tools were projectile points and ulus. These tools were often purpose-made by creating a blank, either by chipping or using a technique where the slate was sawed partway through on one or both sides and then snapped into a blank, then finished by grinding with abraders or whetstones.
Uses
When making the head of an axe out of stone, the piece would be made so it could be hafted. In order to have the stone hafted onto a larger piece, like wood or bone, the ground stone may have at least two notches ground out of one side of the stone, making grooves for the hafting material to lie inside. These grooves would ensure that the stone would not move when struck with a large force. Tough hide would then be wound around the handle and inside the grooves, binding the ground stone and the handle together.
Ground stones were often used as dinnerware. Using large stones, lithic reduction would be done for long periods of time to create bowls and pots for food. Jewelry, beads, ear spools and other decorative ground stones were a sign of high status due to the time and effort needed to make pieces of such small size and detail.
When mashing up seeds and leaves into powders, rounded and smooth ground stones would be used inside a stone bowl. This pair of tools is called a mortar and pestle. The material would be placed into the mortar and the pestle would be moved and pressed into the mortar to grind the material into a fine powder. This process could be used for medicine and cooking. The mortar and pestle are still used today for many cooking recipes.
Ground stones life
Once shaped into an active (making the action, e.g. handstones) or passive (receiving the action, e.g. slab) tool, the stone is used to process a large range of products. At the end of its life, the stone can be recycled - either reused as a ground stone, used as a construction stone, or discarded
Analytical methods
In archaeology, ground stones are studied to document past activities, agriculture strategies and the process of domestication. Ground stones technology was identified as daily life technology associated with cereals processing. However, their study has shown that they are more complex to understand and have multi-function. Ethnographic studies are essential to define the research question but are not enough to correctly determine the function of ground stones. Typology, ethnographic studies and experimental replicas help understand groundstone technologies.
Typical methods are used to analyse them: the trace of wear and the analysis of botanical residues (starch grains and phytoliths).
Use-wear analysis
This technique analyses the trace of wear left by the types of actions (pounding, grinding, cutting). The trace of wear is specific to the action but also the type of product (organic or inorganic) process. Handstones for cereals and hide processing were distinguished using this method.
Botanical remains
Starch grains and phytoliths are botanical remains which are well preserved in archaeological context and are small enough to be trapped in the micro-crevice of a ground stone. Analyzing the botanical remains is essential to precise the function of the stones.
A standardized protocol is used to study plant remains. First, there are extracts from the ground stones, using a toothbrush and an ultrasonic bath. Then, the residues are observed under a polarised microscope. And finish, starch are compared to a reference collection of experimental modern starch grains.
| Technology | Hand tools | null |
147203 | https://en.wikipedia.org/wiki/Polyethylene%20glycol | Polyethylene glycol | {{Chembox
| Verifiedfields = changed
| Watchedfields = changed
| verifiedrevid = 477163023
| Name =
| ImageFile = PEG Structural Formula V1.svg
| ImageFile2 = Polyethylene glycol 400.jpg
| ImageName2 = Polyethylene glycol 400
| IUPACName = poly(oxyethylene) {structure-based}, poly(ethylene oxide) {source-based}
| OtherNames = Kollisolv, Carbowax, GoLYTELY, GlycoLax, Fortrans, TriLyte, Colyte, Halflytely, macrogol, MiraLAX, MoviPrep
| SystematicName =
| Section1 =
| Section2 = {{Chembox Properties
| Formula = C2nH4n+2On+1
| MolarMass = {{nowrap|44.05n + 18.02 g/mol}}
| Appearance =
| Density = 1.125
| MeltingPt =
| BoilingPt =
| Solubility =
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| Section3 =
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Polyethylene glycol (PEG; ) is a polyether compound derived from petroleum with many applications, from industrial manufacturing to medicine. PEG is also known as polyethylene oxide (PEO) or polyoxyethylene (POE), depending on its molecular weight. The structure of PEG is commonly expressed as H−(O−CH2−CH2)n−OH.
Uses
Medical uses
Pharmaceutical-grade PEG is used as an excipient in many pharmaceutical products, in oral, topical, and parenteral dosage forms.
PEG is the basis of a number of laxatives (as MiraLax, RestoraLAX, MoviPrep, etc.). Whole bowel irrigation with polyethylene glycol and added electrolytes is used for bowel preparation before surgery or colonoscopy or for children with constipation. Macrogol (with brand names such as Laxido, Movicol and Miralax) is the generic name for polyethylene glycol used as a laxative. The name may be followed by a number that represents the average molecular weight (e.g. macrogol 3350, macrogol 4000, or macrogol 6000).
The possibility that PEG could be used to fuse axons is being explored by researchers studying peripheral nerve and spinal cord injury.
An example of PEG hydrogels (see Biological uses section) in a therapeutic has been theorized by Ma et al. They propose using the hydrogel to address periodontitis (gum disease) by encapsulating stem cells in the gel that promote healing in the gums. The gel with encapsulated stem cells was to be injected into the site of disease and crosslinked to create the microenvironment required for the stem cells to function.
PEGylation of adenoviruses for gene therapy can help prevent adverse reactions due to pre-existing adenovirus immunity.
A PEGylated lipid is used as an excipient in both the Moderna and Pfizer–BioNTech vaccines for SARS-CoV-2. Both RNA vaccines consist of messenger RNA, or mRNA, encased in a bubble of oily molecules called lipids. Proprietary lipid technology is used for each. In both vaccines, the bubbles are coated with a stabilizing molecule of polyethylene glycol. PEG could trigger allergic reaction, and allergic reactions are the driver for both the United Kingdom and Canadian regulators to issue an advisory, noting that: two "individuals in the U.K. ... were treated and have recovered" from anaphylactic shock. The US CDC stated that in their jurisdiction six cases of "severe allergic reaction" had been recorded from more than 250,000 vaccinations, and of those six only one person had a "history of vaccination reactions".
Chemical uses
Polyethylene glycol is also commonly used as a polar stationary phase for gas chromatography, as well as a heat transfer fluid in electronic testers.
PEG is frequently used to preserve waterlogged wood and other organic artifacts that have been salvaged from underwater archaeological contexts, as was the case with the warship Vasa in Stockholm, and similar cases. It replaces water in wooden objects, making the wood dimensionally stable and preventing the warping or shrinking of the wood when it dries. In addition, PEG is used when working with green wood as a stabilizer, and to prevent shrinkage.
PEG has been used to preserve the painted colors on Terracotta Warriors unearthed at a UNESCO World Heritage site in China. These painted artifacts were created during the Qin Shi Huang (first emperor of China) era. Within 15 seconds of the terra-cotta pieces being unearthed during excavations, the lacquer beneath the paint begins to curl after being exposed to the dry Xi'an air. The paint would subsequently flake off in about four minutes. The German Bavarian State Conservation Office developed a PEG preservative that when immediately applied to unearthed artifacts has aided in preserving the colors painted on the pieces of clay soldiers.
PEG is often used (as an internal calibration compound) in mass spectrometry experiments, with its characteristic fragmentation pattern allowing accurate and reproducible tuning.
PEG derivatives, such as narrow range ethoxylates, are used as surfactants.
PEG has been used as the hydrophilic block of amphiphilic block copolymers used to create some polymersomes.
PEG is a component of the propellent used in UGM-133M Trident II Missiles, in service with the United States Navy.
PEG has been used as a solvent for aryl thioether synthesis.
Biological uses
An example study was done using PEG-diacrylate hydrogels to recreate vascular environments with the encapsulation of endothelial cells and macrophages. This model furthered vascular disease modeling and isolated macrophage phenotype's effect on blood vessels.
PEG is commonly used as a crowding agent in in vitro assays to mimic highly crowded cellular conditions. Although polyethylene glycol is considered biologically inert, it can form non-covalent complexes with monovalent cations such as Na+, K+, Rb+, and Cs+, affecting equilibrium constants of biochemical reactions.
PEG is commonly used as a precipitant for plasmid DNA isolation and protein crystallization. X-ray diffraction of protein crystals can reveal the atomic structure of the proteins.
PEG is used to fuse two different types of cells, most often B-cells and myelomas to create hybridomas. César Milstein and Georges J. F. Köhler originated this technique, which they used for antibody production, winning a Nobel Prize in Physiology or Medicine in 1984.
In microbiology, PEG precipitation is used to concentrate viruses. PEG is also used to induce complete fusion (mixing of both inner and outer leaflets) in liposomes reconstituted in vitro.
Gene therapy vectors (such as viruses) can be PEG-coated to shield them from inactivation by the immune system and to de-target them from organs where they may build up and have a toxic effect. The size of the PEG polymer is important, with larger polymers achieving the best immune protection.
PEG is a component of stable nucleic acid lipid particles (SNALPs) used to package siRNA for use in vivo. (free with registration)
In blood banking, PEG is used as a potentiator to enhance detection of antigens and antibodies.
When working with phenol in a laboratory situation, PEG 300 can be used on phenol skin burns to deactivate any residual phenol.
In biophysics, polyethylene glycols are the molecules of choice for the functioning ion channel diameter studies, because in aqueous solutions they have a spherical shape and can block ion channel conductance.
Commercial uses
PEG is the basis of many skin creams (as cetomacrogol) and personal lubricants.
PEG is used in a number of toothpastes as a dispersant. In this application, it binds water and helps keep xanthan gum uniformly distributed throughout the toothpaste.
PEG is under investigation for use in liquid body armor, and in tattoos to monitor diabetes.
Polymer segments derived from PEG polyols impart flexibility to polyurethanes for applications such as elastomeric fibers (spandex) and foam cushions.
In low-molecular-weight formulations (e.g. PEG 400), it is used in Hewlett-Packard designjet printers as an ink solvent and lubricant for the print heads.
PEG is used as an anti-foaming agent in food and drinks – its INS number is 1521 or E1521 in the EU.
Industrial uses
A nitrate ester-plasticized polyethylene glycol (NEPE-75) is used in Trident II submarine-launched ballistic missile solid rocket fuel.
Dimethyl ethers of PEG are the key ingredient of Selexol, a solvent used by coal-burning, integrated gasification combined cycle (IGCC) power plants to remove carbon dioxide and hydrogen sulfide from the syngas stream.
PEG has been used as the gate insulator in an electric double-layer transistor to induce superconductivity in an insulator.
PEG is used as a polymer host for solid polymer electrolytes. Although not yet in commercial production, many groups around the globe are engaged in research on solid polymer electrolytes involving PEG, to improve their properties, and in permitting their use in batteries, electro-chromic display systems, and other products in the future.
PEG is injected into industrial processes to reduce foaming in separation equipment.
PEG is used as a binder in the preparation of technical ceramics.
PEG was used as an additive to silver halide photographic emulsions.
Entertainment uses
PEG is used to extend the size and durability of very large soap bubbles.
PEG is an ingredient in some personal lubricants. (Not to be confused with propylene glycol.)
PEG is the main ingredient in the paint (known as "fill") in paintballs.
Human health effects
PEO's have "very low single dose oral toxicity", on the order of tens of grams per kilogram of human body weight when ingested by mouth. Because of its low toxicity, PEO is used in a variety of edible products. It is also used as a lubricating coating for various surfaces in aqueous and non-aqueous applications.
The precursor to PEGs is ethylene oxide, which is hazardous. Ethylene glycol and its ethers are nephrotoxic (poisonous to the kidneys) if applied to damaged skin.
The United States Food and Drug Administration (FDA or US FDA) regards PEG as biologically inert and safe.
A 2015 study appears to challenge the FDA's conclusion. In the study, a high-sensitivity ELISA assay detected anti-PEG antibodies in 72% of random blood plasma samples collected from 1990 to 1999. According to the study's authors, this result suggests that anti-PEG antibodies may be present, typically at low levels, in people who were never treated with PEGylated drugs. Due to its ubiquity in many products and the large percentage of the population with antibodies to PEG, which indicates an allergic reaction, hypersensitive reactions to PEG are an increasing health concern. Allergy to PEG is usually discovered after a person has been diagnosed with an allergy to several seemingly unrelated products—including processed foods, cosmetics, drugs, and other substances—that contain or were manufactured with PEG.
Available forms and nomenclaturePEG, PEO, and POE refer to an oligomer or polymer of ethylene oxide. The three names are chemically synonymous, but historically PEG is preferred in the biomedical field, whereas PEO is more prevalent in the field of polymer chemistry. Because different applications require different polymer chain lengths, PEG has tended to refer to oligomers and polymers with a molecular mass below 20,000g/mol, PEO to polymers with a molecular mass above 20,000g/mol, and POE to a polymer of any molecular mass. PEGs are prepared by polymerization of ethylene oxide and are commercially available over a wide range of molecular weights from 300g/mol to 10,000,000g/mol.
PEG and PEO are liquids or low-melting solids, depending on their molecular weights. While PEG and PEO with different molecular weights find use in different applications and have different physical properties (e.g. viscosity) due to chain length effects, their chemical properties are nearly identical. Different forms of PEG are also available, depending on the initiator used for the polymerization process – the most common initiator is a monofunctional methyl ether PEG, or methoxypoly(ethylene glycol), abbreviated mPEG. Lower-molecular-weight PEGs are also available as purer oligomers, referred to as monodisperse, uniform, or discrete. Very high-purity PEG has recently been shown to be crystalline, allowing the determination of a crystal structure by x-ray crystallography. Since purification and separation of pure oligomers is difficult, the price for this type of quality is often 10–1000 fold that of polydisperse PEG.
PEGs are also available with different geometries.
Branched PEGs have three to ten PEG chains emanating from a central core group.
Star PEGs have 10 to 100 PEG chains emanating from a central core group.
Comb PEGs have multiple PEG chains normally grafted onto a polymer backbone.
The numbers that are often included in the names of PEGs indicate their average molecular weights (e.g. a PEG with would have an average molecular weight of approximately 400 daltons, and would be labeled PEG 400). Most PEGs include molecules with a distribution of molecular weights (i.e. they are polydisperse). The size distribution can be characterized statistically by its weight average molecular weight (Mw) and its number average molecular weight (Mn), the ratio of which is called the polydispersity index (ĐM). Mw and Mn can be measured by mass spectrometry.
PEGylation is the act of covalently coupling a PEG structure to another larger molecule, for example, a therapeutic protein, which is then referred to as a PEGylated protein. PEGylated interferon alfa-2a or alfa-2b are commonly used injectable treatments for hepatitis C infection.
PEG is soluble in water, methanol, ethanol, acetonitrile, benzene, and dichloromethane, and is insoluble in diethyl ether and hexane. It is coupled to hydrophobic molecules to produce non-ionic surfactants.
PEG and related polymers (PEG phospholipid constructs) are often sonicated when used in biomedical applications. However, as reported by Murali et al., PEG is very sensitive to sonolytic degradation and PEG degradation products can be toxic to mammalian cells. It is, thus, imperative to assess potential PEG degradation to ensure that the final material does not contain undocumented contaminants that can introduce artifacts into experimental results.
PEGs and methoxypolyethylene glycols are manufactured by Dow Chemical under the trade name Carbowax for industrial use, and Carbowax Sentry'' for food and pharmaceutical use. They vary in consistency from liquid to solid, depending on the molecular weight, as indicated by a number following the name. They are used commercially in numerous applications, including foods, cosmetics, pharmaceutics, biomedicine, dispersing agents, solvents, ointments, suppository bases, as tablet excipients, and as laxatives. Some specific groups are lauromacrogols, nonoxynols, octoxynols, and poloxamers.
Production
The production of polyethylene glycol was first reported in 1859. Both A. V. Lourenço and Charles Adolphe Wurtz independently isolated products that were polyethylene glycols. Polyethylene glycol is produced by the interaction of ethylene oxide with water, ethylene glycol, or ethylene glycol oligomers. The reaction is catalyzed by acidic or basic catalysts. Ethylene glycol and its oligomers are preferable as a starting material instead of water because they allow the creation of polymers with a low polydispersity (narrow molecular weight distribution). Polymer chain length depends on the ratio of reactants.
HOCH2CH2OH + n(CH2CH2O) → HO(CH2CH2O)n+1H
Depending on the catalyst type, the mechanism of polymerization can be cationic or anionic. The anionic mechanism is preferable because it allows one to obtain PEG with a low polydispersity. Polymerization of ethylene oxide is an exothermic process. Overheating or contaminating ethylene oxide with catalysts, such as alkalis or metal oxides, can lead to runaway polymerization, which can end in an explosion after a few hours.
Polyethylene oxide, or high-molecular-weight polyethylene glycol, is synthesized by suspension polymerization. It is necessary to hold the growing polymer chain in solution in the course of the polycondensation process. The reaction is catalyzed by magnesium-, aluminium-, or calcium-organoelement compounds. To prevent coagulation of polymer chains from solution, chelating additives, such as dimethylglyoxime, are used.
Alkaline catalysts, such as sodium hydroxide (NaOH), potassium hydroxide (KOH), or sodium carbonate (Na2CO3), are used to prepare low-molecular-weight polyethylene glycol.
| Physical sciences | Polymers | Chemistry |
147235 | https://en.wikipedia.org/wiki/Chinese%20units%20of%20measurement | Chinese units of measurement | Chinese units of measurement, known in Chinese as the shìzhì ("market system"), are the traditional units of measurement of the Han Chinese. Although Chinese numerals have been decimal (base-10) since the Shang, several Chinese measures use hexadecimal (base-16). Local applications have varied, but the Chinese dynasties usually proclaimed standard measurements and recorded their predecessor's systems in their histories.
In the present day, the People's Republic of China maintains some customary units based upon the market units but standardized to round values in the metric system, for example the common jin or catty of exactly 500g. The Chinese name for most metric units is based on that of the closest traditional unit; when confusion might arise, the word "market" (, shì) is used to specify the traditional unit and "common" or "public" (, gōng) is used for the metric value. Taiwan, like Korea, saw its traditional units standardized to Japanese values and their conversion to a metric basis, such as the Taiwanese ping of about 3.306m2 based on the square ken. The Hong Kong SAR continues to use its traditional units, now legally defined based on a local equation with metric units. For instance, the Hong Kong catty is precisely .
Note: The names lí ( or ) and fēn () for small units are the same for length, area, and mass; however, they refer to different kinds of measurements.
History
According to the Liji, the legendary Yellow Emperor created the first measurement units. The Xiao Erya and the Kongzi Jiayu state that length units were derived from the human body. According to the Records of the Grand Historian, these human body units caused inconsistency, and Yu the Great, another legendary figure, unified the length measurements. Rulers with decimal units have been unearthed from Shang dynasty tombs.
In the Zhou dynasty, the king conferred nobles with powers of the state and the measurement units began to be inconsistent from state to state. After the Warring States period, Qin Shi Huang unified China, and later standardized measurement units. In the Han dynasty, these measurements were still being used, and were documented systematically in the Book of Han.
Astronomical instruments show little change of the length of chi in the following centuries, since the calendar needed to be consistent. It was not until the introduction of decimal units in the Ming dynasty that the traditional system was revised.
Republican Era
On 7 January 1915, the Beiyang government promulgated a measurement law to use not only metric system as the standard but also a set of Chinese-style measurement based directly on the Qing dynasty definitions ().
On 16 February 1929, the Nationalist government adopted and promulgated The Weights and Measures Act to adopt the metric system as the official standard and to limit the newer Chinese units of measurement () to private sales and trade in Article 11, effective on 1 January 1930. These newer "market" units are based on rounded metric numbers.
People's Republic of China
The Government of the People's Republic of China continued using the market system along with metric system, as decreed by the State Council of the People's Republic of China on 25 June 1959, but 1 catty being 500 grams, would become divided into 10 (new) taels, instead of 16 (old) taels, to be converted from province to province, while exempting Chinese prescription drugs from the conversion to prevent errors.
On 27 February 1984, the State Council of the People's Republic of China decreed the market system to remain acceptable until the end of 1990 and ordered the transition to the national legal measures by that time, but farmland measures would be exempt from this mandatory metrication until further investigation and study.
Hong Kong
In 1976 the Hong Kong Metrication Ordinance allowed a gradual replacement of the system in favor of the International System of Units (SI) metric system. The Weights and Measures Ordinance defines the metric, Imperial, and Chinese units. As of 2012, all three systems are legal for trade and are in widespread use.
Macau
On 24 August 1992, Macau published Law No. 14/92/M to order that Chinese units of measurement similar to those used in Hong Kong, Imperial units, and United States customary units would be permissible for five years since the effective date of the Law, 1 January 1993, on the condition of indicating the corresponding SI values, then for three more years thereafter, Chinese, Imperial, and US units would be permissible as secondary to the SI.
Ancient Chinese units
Length
Traditional units of length include the chi (), bu (), and li (). The precise length of these units, and the ratios between these units, has varied over time. 1 bu has consisted of either 5 or 6 chi, while 1 li has consisted of 300 or 360 bu.
Modern Chinese units
All "metric values" given in the tables are exact unless otherwise specified by the approximation sign '~'.
Certain units are also listed at List of Chinese classifiers → Measurement units.
Length
Chinese length units promulgated in 1915
Chinese length units effective in 1930
Metric length units
The Chinese word for metre is mǐ; this can take the Chinese standard SI prefixes (for "kilo-", "centi-", etc.). A kilometre, however, may also be called gōnglǐ, i.e. a metric lǐ.
In the engineering field, traditional units are rounded up to metric units. For example, the Chinese word (T) or (S) sī is used to express 0.01 mm.
Hong Kong and Macau length units
These correspond to the measures listed simply as "China" in The Measures, Weights, & Moneys of All Nations
Area
Chinese area units promulgated in 1915
Chinese area units effective in 1930
Metric and other area units
Metric and other standard length units can be squared by the addition of the prefix píngfāng. For example, a square kilometre is píngfāng gōnglǐ.
Macau area units
Volume
These units are used to measure cereal grains, among other things. In imperial times, the physical standard for these was the jialiang.
Chinese volume units promulgated in 1915
Chinese volume units effective in 1930
Metric volume units
In the case of volume, the market and metric shēng coincide, being equal to one litre as shown in the table. The Chinese standard SI prefixes (for "milli-", "centi-", etc.) may be added to this word shēng.
Units of volume can also be obtained from any standard unit of length using the prefix lìfāng ("cubic"), as in lìfāng mǐ for one cubic metre.
Macau volume units
Mass
These units are used to measure the mass of objects. They are also famous for measuring monetary objects such as gold and silver.
Chinese mass units promulgated in 1915
Mass units in the Republic of China since 1930
Mass units in the People's Republic of China since 1959
Metric mass units
The Chinese word for gram is kè; this can take the Chinese standard SI prefixes (for "milli-", "deca-", and so on). A kilogram, however, is commonly called gōngjīn, i.e. a metric jīn.
Hong Kong and Macau mass units
Hong Kong troy units
These are used for trading precious metals such as gold and silver.
Time
Historiography
As there were hundreds of unofficial measures in use, the bibliography is quite vast. The editions of Wu Chenglou's 1937 History of Chinese Measurement were the usual standard up to the 1980s or so, but rely mostly on surviving literary accounts. Newer research has put more emphasis on archeological discoveries. Qiu Guangming & Zhang Yanming's 2005 bilingual Concise History of Ancient Chinese Measures and Weights summarizes these findings. A relatively recent and comprehensive bibliography, organized by period studied, has been compiled in 2012 by Cao & al.; for a shorter list, see Wilkinson's year 2000 Chinese History.
| Physical sciences | Measurement systems | Basics and measurement |
147242 | https://en.wikipedia.org/wiki/Octane%20rating | Octane rating | An octane rating, or octane number, is a standard measure of a fuel's ability to withstand compression in an internal combustion engine without causing engine knocking. The higher the octane number, the more compression the fuel can withstand before detonating. Octane rating does not relate directly to the power output or the energy content of the fuel per unit mass or volume, but simply indicates the resistance to detonating under pressure without a spark.
Whether a higher octane fuel improves or impairs an engine's performance depends on the design of the engine. In broad terms, fuels with a higher octane rating are used in higher-compression gasoline engines, which may yield higher power for these engines. The added power in such cases comes from the way the engine is designed to compress the air/fuel mixture, and not directly from the rating of the gasoline.
In contrast, fuels with lower octane (but higher cetane numbers) are ideal for diesel engines because diesel engines (also called compression-ignition engines) do not compress the fuel, but rather compress only air, and then inject fuel into the air that was heated by compression. Gasoline engines rely on ignition of compressed air and fuel mixture, which is ignited only near the end of the compression stroke by electric spark plugs. Therefore, being able to compress the air/fuel mixture without causing detonation is important mainly for gasoline engines. Using gasoline with lower octane than an engine is built for may cause engine knocking and/or pre-ignition.
The octane rating of aviation gasoline was extremely important in determining aero engine performance in the aircraft of World War II. The octane rating affected not only the performance of the gasoline, but also its versatility; the higher octane fuel allowed a wider range of lean to rich operating conditions.
Principles
The problem: detonation
In spark ignition internal combustion engines, knocking (also knock, detonation, spark knock, pinging, or pinking) occurs when combustion of some of the air/fuel mixture in the cylinder does not result from propagation of the flame front ignited by the spark plug, but when one or more pockets of air/fuel mixture explode outside the envelope of the normal combustion front. The fuel-air charge is meant to be ignited by the spark plug only, and at a precise point in the piston's stroke. Knock occurs when the peak of the combustion process no longer occurs at the optimum moment for the four-stroke cycle. In a simple explanation, the forward moving wave of combustion that burns the hydrocarbon + oxygen mixture inside the cylinder like a wave that a surfer would wish to surf upon is violently disrupted by a secondary wave that has started elsewhere. The shock wave of these two separate waves creates the characteristic metallic "pinging" sound, and cylinder pressure increases dramatically. Effects of engine knocking range from inconsequential (incremental heating plus power loss) to completely destructive (detonation while one of the valves is still open).
Knocking should not be confused with pre-ignitionthey are two separate events with pre-ignition occurring before the combustion event. However, pre-ignition is highly correlated with knock because knock will cause rapid heat increase within the cylinder eventually leading to destructive pre-detonation.
Most engine management systems commonly found in automobiles today, typically electronic fuel injection (EFI), have a knock sensor that monitors if knock is being produced by the fuel being used. In modern computer-controlled engines, the ignition timing will be automatically altered by the engine management system to reduce the knock to an acceptable level.
Iso-octane as a reference standard
Octanes are a family of hydrocarbons that are typical components of gasoline. They are colorless liquids that boil around 125 °C (260 °F). One member of the octane family, 2,2,4-Trimethylpentane (iso-octane), is used as a reference standard to benchmark the tendency of gasoline or LPG fuels to resist self-ignition.
The octane rating of gasoline is measured in a test engine and is defined by comparison with the mixture of 2,2,4-trimethylpentane (iso-octane) and normal heptane that would have the same anti-knocking capability as the fuel under test. The percentage, by volume, of 2,2,4-trimethylpentane in that mixture is the octane number of the fuel. For example, gasoline with the same knocking characteristics as a mixture of 90% iso-octane and 10% heptane would have an octane rating of 90. A rating of 90 does not mean that the gasoline contains just iso-octane and heptane in these proportions, but that it has the same detonation resistance properties (generally, gasoline sold for common use never consists solely of iso-octane and heptane; it is a mixture of many hydrocarbons and often other additives).
Octane ratings are not indicators of the energy content of fuels. (See Effects below and Heat of combustion). They are only a measure of the fuel's tendency to burn in a controlled manner, rather than exploding in an uncontrolled manner.
Where the octane number is raised by blending in ethanol, energy content per volume is reduced. Ethanol energy density can be compared with gasoline in heat-of-combustion tables.
It is possible for a fuel to have a Research Octane Number (RON) more than 100, because iso-octane is not the most knock-resistant substance available today. Racing fuels, avgas, LPG and alcohol fuels such as methanol may have octane ratings of 110 or significantly higher. Typical "octane booster" gasoline additives include MTBE, ETBE, iso-octane and toluene. Lead in the form of tetraethyllead was once a common additive, but concerns about its toxicity have led to its use for fuels for road vehicles being progressively phased out worldwide beginning in the 1970s.
Measurement methods
Research Octane Number (RON)
The most common type of octane rating worldwide is the Research Octane Number (RON). RON is determined by running the fuel in a test engine at 600 rpm with a variable compression ratio under controlled conditions, and comparing the results with those for mixtures of iso-octane and n-heptane. The compression ratio is varied during the test to challenge the fuel's antiknocking tendency, as an increase in the compression ratio will increase the chances of knocking.
Motor Octane Number (MON)
Another type of octane rating, called Motor Octane Number (MON), is determined at 900 rpm engine speed instead of the 600 rpm for RON. MON testing uses a similar test engine to that used in RON testing, but with a preheated fuel mixture, higher engine speed, and variable ignition timing to further stress the fuel's knock resistance. Depending on the composition of the fuel, the MON of a modern pump gasoline will be about 8 to 12 lower than the RON, but there is no direct link between RON and MON. See the table below.
Anti-Knock Index (AKI) or (R+M)/2
In Canada, The United States, and Mexico, the advertised octane rating is the average of the RON and the MON, called the Anti-Knock Index (AKI). It is often written on pumps as (R+M)/2. AKI is also sometimes called PON (Pump Octane Number).
Difference between RON, MON, and AKI
Because of the 8 to 12 octane number difference between RON and MON noted above, the AKI shown in Canada and the United States is 4 to 6 octane numbers lower than elsewhere in the world for the same fuel. This difference between RON and MON is known as the fuel's sensitivity, and is not typically published for those countries that use the Anti-Knock Index labelling system.
See the table in the following section for a comparison.
Observed Road Octane Number (RdON)
Another type of octane rating, called Observed Road Octane Number (RdON), is derived from testing the gasoline in ordinary multi-cylinder engines (rather than in a purpose-built test engine), normally at wide open throttle. This type of test was developed in the 1920s and is still reliable today. The original RdON tests were done in cars on the road, but as technology developed the testing was moved to chassis dynamometers with environmental controls to improve consistency.
Octane Index
The evaluation of the octane number by either of the two laboratory methods requires a special engine built to match the tests' rigid standards, and the procedure can be both expensive and time-consuming. The standard engine required for the test may not always be available, especially in out-of-the-way places or in small or mobile laboratories. These and other considerations led to the search for a rapid method for the evaluation of the anti-knock quality of gasoline. Such substitute methods include FTIR, near infrared on-line analyzers, and others. Deriving an equation that can be used to calculate ratings accurately enough would also serve the same purpose, with added advantages. The term Octane Index is often used to refer to the use of an equation to determine a theoretical rating, in contradistinction to the direct measurements required for research or motor octane numbers. An octane index can be of great service in the blending of gasoline. Motor gasoline, as marketed, is usually a blend of several types of refinery grades that are derived from different processes such as straight-run gasoline, reformate, cracked gasoline etc. These different grades are blended in amounts that will meet final product specifications. Most refiners produce and market more than one grade of motor gasoline, differing principally in their anti-knock quality. Being able to make sufficiently accurate estimates of the octane rating that will result from blending different refinery products is essential, something for which the calculated octane index is specially suited.
Aviation gasoline octane ratings
Aviation gasolines used in piston aircraft engines common in general aviation have a slightly different method of measuring the octane of the fuel. Similar to an AKI, it has two different ratings, although it is usually referred to only by the lower of the two. One is referred to as the "aviation lean" rating, which for ratings up to 100 is the same as the MON of the fuel. The second is the "aviation rich" rating and corresponds to the octane rating of a test engine under forced induction operation common in high-performance and military piston aircraft. This utilizes a supercharger, and uses a significantly richer fuel/air ratio for improved detonation resistance.
The most common currently used fuel, 100LL, has an aviation lean rating of 100 octane, and an aviation rich rating of 130.
Examples
The RON/MON values of n-heptane and iso-octane are exactly 0 and 100, respectively, by the definition of octane rating. The following table lists octane ratings for various other fuels.
Effects
Higher octane ratings correlate to higher activation energies: the amount of applied energy required to initiate combustion. Since higher octane fuels have higher activation energy requirements, it is less likely that a given compression will cause uncontrolled ignition, otherwise known as autoignition, self-ignition, pre-ignition, detonation, or knocking.
Because octane is a measured and/or calculated rating of the fuel's ability to resist autoignition, the higher the octane of the fuel, the harder that fuel is to ignite and the more heat is required to ignite it. The result is that a hotter ignition spark is required for ignition. Creating a hotter spark requires more energy from the ignition system, which in turn increases the parasitic electrical load on the engine. The spark also must begin earlier in order to generate sufficient heat at the proper time for precise ignition. As octane, ignition spark energy, and the need for precise timing increase, the engine becomes more difficult to "tune" and keep "in tune". The resulting sub-optimal spark energy and timing can cause major engine problems, from a simple "miss" to uncontrolled detonation and catastrophic engine failure.
Mechanically within the cylinder, stability can be visualized as having a flame wave initiate at the spark plug and then "travel in a fairly uniform manner across the combustion chamber" with the expanding gas mix pushing the piston throughout the entirety of the power stroke. A stable gasoline and air mix will combust when the flame wave reaches the molecules, adding heat at the interface. Knock occurs when a secondary flame wave forms from instability and then travels against the path of the primary flame wave, thus depriving the power stroke of its uniformity and causing issues including power loss and heat buildup.
The other rarely-discussed reality with high-octane fuels associated with "high performance" is that as octane increases, the specific gravity and energy content of the fuel per unit of weight are reduced. The net result is that to make a given amount of power, more high-octane fuel must be burned in the engine. Lighter and "thinner" fuel also has a lower specific heat, so the practice of running an engine "rich" to use excess fuel to aid in cooling requires richer and richer mixtures as octane increases.
Higher-octane, lower-energy-dense "thinner" fuels often contain alcohol compounds incompatible with the stock fuel system components, which also makes them hygroscopic. They also evaporate away much more easily than heavier, lower-octane fuel which leads to more accumulated contaminants in the fuel system. It is typically the and the compounds in the fuel that have the most detrimental effects on the engine fuel system components, as such acids corrode many metals used in gasoline fuel systems.
During the compression stroke of an internal combustion engine, the temperature of the air-fuel mix rises as it is compressed, in accordance with the ideal gas law. Higher compression ratios necessarily add parasitic load to the engine, and are only necessary if the engine is being specifically designed to run on high-octane fuel. Aircraft engines run at relatively low speeds and are "undersquare". They run best on lower-octane, slower-burning fuels that require less heat and a lower compression ratio for optimum vaporization and uniform fuel-air mixing, with the ignition spark coming as late as possible in order to extend the production of cylinder pressure and torque as far down the power stroke as possible. The main reason for using high-octane fuel in air-cooled engines is that it is more easily vaporized in a cold carburetor and engine and absorbs less intake air heat which greatly reduces the tendency for carburetor icing to occur.
With their reduced densities and weight per volume of fuel, the other obvious benefit is that an aircraft with any given volume of fuel in the tanks is automatically lighter. And since many airplanes are flown only occasionally and may sit unused for weeks or months, the lighter fuels tend to evaporate away and leave behind fewer deposits such as "varnish" (gasoline components, particularly alkenes and oxygenates slowly polymerize into solids). Aircraft also typically have dual "redundant" ignition systems which are nearly impossible to tune and time to produce identical ignition timing, so using a lighter fuel that's less prone to autoignition is a wise "insurance policy". For the same reasons, those lighter fuels which are better solvents are much less likely to cause any "varnish" or other fouling on the "backup" spark plugs.
In almost all general aviation piston engines, the fuel mixture is directly controlled by the pilot, via a knob and cable or lever similar to (and next to) the throttle control. Leaningreducing the mixture from its maximum amount – must be done with knowledge, as some combinations of fuel mixture and throttle position (that produce the highest ) can cause detonation and/or pre-ignition, in the worst case destroying the engine within seconds. Pilots are taught in primary training to avoid settings that produce the highest exhaust gas temperatures, and run the engine either "rich of peak EGT" (more fuel than can be burned with the available air) or "lean of peak" (less fuel, leaving some oxygen in the exhaust) as either will keep the fuel-air mixture from detonating prematurely. Because of the high cost of unleaded, high-octane avgas, and possible increased range before refueling, some general aviation pilots attempt to save money by tuning their fuel-air mixtures and ignition timing to run "lean of peak". Additionally, the decreased air density at higher altitudes (such as Colorado) and temperatures (as in summer) requires leaning (reduction in amount of fuel per volume or mass of air) for the peak EGT and power (crucial for takeoff).
Regional variations
The selection of octane ratings available at filling stations can vary greatly between countries.
Australia: "regular" unleaded fuel is 91 RON, "premium" unleaded with 95 RON is widely available, and 98 RON fuel is also very common. Shell used to sell 100 RON fuel (5% ethanol content) from a small number of service stations, most of which are located in major cities (stopped in August 2008). United Petroleum used to sell 100 RON unleaded fuel (10% ethanol content) at a small number of its service stations (originally only two, but then expanded to 67 outlets nationwide) (stopped in September 2014). All fuel in Australia is unleaded except for some aviation fuels. E85 unleaded fuel is also available at several United service stations across the country. By 2018, E10 fuel had become quite common, and is available at almost every major fuel station, except in Western Australia.
Bahrain: 91 and 95 (RON), standard in all gasoline stations in the country and advertised as (Jayyid) for Regular or 91 and (Mumtaz) for Premium or 95 and 98 (RON) as super.
Bangladesh: Two types of fuel are available at petrol stations in Bangladesh. Motor Gasoline Regular (marketed as "Petrol") which has RON 80 rating, and Motor Gasoline Premium (marketed as "Octane") which has RON 95 rating. Petrol stations in Bangladesh are privatised, but the prices are regulated by the authorities and have a fixed price at BDT 86.00 (US$1.04) and BDT 89.00 (US$1.07) (as of 1 March 2018) per litre respectively.
Botswana: 93 and 95 RON are standard at almost all gas stations thorough Botswana. The two types are unleaded.
Brazil: As defined by federal law, the RON standard is used and all types of gasoline sold in all gas stations throughout the country are unleaded (the latter since 1991). By default, it was defined by the federal government that the regular (and the lowest) octane standard in Brazil is 93 RON, known in Portuguese as Gasolina Comum (English: "Common Gasoline") – Petrobras stations brand it as Gasolina Regular (English: "Regular Gasoline"). This type of gasoline can be found in most Brazilian petrol stations and does not have any additives, except the inclusion of 27,5% of ethanol (as required by the Brazilian National Agency of Petroleum, Natural Gas and Biofuels – Portuguese: Agência Nacional do Petróleo, Gás Natural e Biocombustíveis or simply ANP – since 2011). Along with the "Common" gasoline, there is a second type of gasoline that can also be found in most stations in Brazil. This gasoline is also mixed with 27,5% of ethanol (to comply with the ANP regulation, that prohibits the sale of the 100% "pure gasoline" compound in all Brazilian stations), but a few detergent and dispersant additives are also included in the compound. This type of gasoline is known in Portuguese as Gasolina Aditivada (English: "Additived Gasoline") – Petrobras stations brand it as "Petrobras Grid"; nevertheless, the octane rating is also 93 RON (these additives are used to improve the performance and efficiency of the engine, but they are not indicative of a higher octane rating). However, higher octane levels of gasoline are found in many stations (all stations in Brazil, regardless of the octane rating, have to conform the ANP requirement of 27,5% of ethanol mixed with the gasoline, and both "Common" and "Additived" gasolines can also be found in most of these stations), such as the "Premium Gasoline" (known in Portuguese as Gasolina Premium – 98 RON), the "OctaPro" (103 RON), sold at Ipiranga stations, and the "Petrobras Podium" (102 RON), sold at Petrobras stations.
Canada: in Canada octane rating is displayed in AKI. In most areas, the standard grades are 87 (regular), 89 (mid-grade) and 91–94 (premium) AKI. In the Atlantic Provinces, gasoline is often available without any blend of ethanol, but only up to 91 AKI.
China: From January 1, 2000, all fuel stations offer unleaded fuel only. Now,92 RON and 95 RON (previously 90 RON, 93 RON and 97 RON) are commonly offered. Some state-run gas stations (Sinopec, PetroChina) in various cities sell 98 RON, but not all. Private gas stations outside of China's Shandong province rarely offer 98 RON. In most rural areas it can be difficult to find fuel with over 95 RON. In backward provinces and regions, only ethanol gasoline containing 10% ethanol is allowed to be sold: 92E10, 95E10 and 98E10, Some gas pumps use the labels "E92, E95 and E98", but they still represent E10 ethanol gasoline of 92 RON, 95 RON and 98 RON. Sinopec's 98 RON gasoline is called "X-power", and PetroChina's 98 RON gasoline is called "CN98". China's National VI gasoline standard has completely banned the use of metal anti-knock agents, because metal anti-knock agents such as MMT and ferrocene will clog the car's GPF, but gasoline vehicles that meet the National VI emission standards must install GPF.
Chile: 93, 95 and 97 RON are standard at almost all gas stations thorough Chile. The three types are unleaded.
Colombia: "Ecopetrol", Colombia's monopoly of refining and distribution of gasoline establishes a minimum AKI of 81 octanes for "Corriente" gasoline and minimum AKI of 87 octanes for "Extra" gasoline. (91.5 RON corriente, and 95 RON for extra)
Costa Rica: RECOPE, Costa Rica's distribution monopoly, establishes the following ratings: Plus 91 (at least 91 RON) and Super (at least 95 RON).
Croatia: All fuel stations offer unleaded "Eurosuper BS" (abbreviation "BS" meaning "no sulfur content") 95 RON fuel, many also offer "Eurosuper Plus BS" 98 RON. Some companies offer 100 RON fuel instead of 98.
Cyprus: All fuel stations offer unleaded 95 and 98 RON, and a few offer 100 RON as well.
Denmark: 95 RON is a common budget choice, with 95 and 98 being widely available, and 92 rarely seen as it has been phased out during the 2010s. A selection of brands offers >=100 options, under trademarked names. However several fuel stations are phasing out 92 RON. By law, it is decided that all gasoline companies from July 2010 to January 2020 should use a mix containing 5% bioethanol in the gasoline and increased to 10% after January 2020.
Ecuador: "Extra" and Ecopais (5% etanol) with 85 RON, "Eco Plus" with 89 RON and "Super Premium" with 95 (RON). Extra/Ecopais and Super Premium are available in all fuel stations. "Extra" is the most commonly used. All fuels are unleaded.
Egypt: Egyptian fuel stations had 90 RON until July 2014 when the government found no remaining use for it, leaving only 92 RON and 95 RON. 80 RON is found in a very limited number of fuel stations as they are used only for extremely old cars that cannot cope with high octane fuel. 95 RON was used limitedly due to its high price (more than twice the price of 92 RON). But after the increasing the prices again in 2018, 95 RON price became only 15% higher than 92 RON, so it started to gain popularity.
Estonia: 95 RON and 98 RON are widely available. E85 (bioethanol) gasoline found in very few gas stations.
Finland: 95 and 98 (RON), advertised as such, at almost all gas stations. Most cars run on 95, but 98 is available for vehicles that need higher octane fuel, or older models containing parts easily damaged by high ethanol content. Shell offers V-Power, advertised as "over 99 octane", instead of 98. In the beginning of 2011 95 RON was replaced by 95E10 containing 10% ethanol, and 98 RON by 98E5, containing 5% ethanol. ST1 also offers RE85 on some stations, which is 85% ethanol made from biodegradable waste (from which the advertised name "ReFuel" comes). RE85 is only suitable for flexifuel cars that can run on high-percentage ethanol.
Germany: "Super E5 and E10" 95 RON and "Super Plus E5" 98 RON are available practically everywhere. Big suppliers such as Shell or Aral offer 100 RON gasoline (Shell V-Power, Aral Ultimate) at almost every fuel station. "Normal" 91 RON is only rarely offered because lower production amounts make it more expensive than "Super" 95 RON. Due to a new European Union law, gas stations are being required to offer a minimum rate of the new mixture of "Super" 95 RON with up to 10% ethanol branded as "Super E10".
Greece (Hellas): 95 RON (standard unleaded), 98 & 100 RON unleaded offered by some companies (e.g., EKO, Shell, BP).
Hong Kong: only 98 RON is available in the market. There have been calls to re-introduce 95 RON, but the calls have been rejected by all automotive fuel station chains, citing that 95 RON was phased out because of market forces.
India: India's ordinary and premium petrol options are of 91 RON. The premium petrols are generally ordinary fuels with additives, that do not really change the octane value. Two variants, "Speed 93" and "Speed 97", were launched, with RON values of 93 and 97. In 2017, Hindustan Petroleum launched poWer 99 with an RON value of 99 which was initially available only in Bangalore, Pune and now in Mumbai but was expected to roll out in other major cities soon. India's economy-class vehicles usually have compression ratios under 10:1, thus enabling them to use lower-octane petrol without engine knocking.
Indonesia: Indonesia's "Premium" gasoline, rated at 88 RON, was the lowest grade gasoline, but was phased out by 2021. Other options have been "Pertalite", rated at 90 RON; "Pertamax", rated at 92 RON; "Pertamax Plus", rated at 95 RON (now replaced by Pertamax Green in July 2023); and "Pertamax Racing", a 100 RON fuel sold in selected stations. From August 2016, Pertamina began selling "Pertamax Turbo", rated at 98 RON, as a replacement for Pertamax Plus. Total and Shell stations only sell RON 92 and 95 gasoline. Shell launched a new variant, "Regular", rated at 90 RON, in early 2018, but this was discontinued in January 2022. However, after 6 years of discontinuation of Pertamax Plus, In July 2023 Pertamina launched the Pertamax Green 95 which made of sugarcane and a mixture of Pertamax 92 and the price is slightly cheaper than Pertamax Turbo which rated 98 RON.
Iran: 'regular' gasoline has an octane rating of 87 RON, which is the most prevalent type of gasoline available throughout the country. Select gas stations within major cities also offer 'Super' 95 RON. Due to high air pollution, an environmentally cleaner variety, marketed as Euro-4, is being introduced in metropolitan areas instead of the Regular, which boasts an octane rating of 91 RON and sulphur levels not exceeding 50 ppm.
Ireland: 95 RON "unleaded" is the only gasoline type available through stations, although E5 (99 RON) is becoming more commonplace.
Italy: 95 RON is the only compulsory gasoline offered (verde, "green"), only a few fuel stations (Agip, IP, IES, OMV) offer 98 RON as the premium type, many Shell and Tamoil stations close to the cities offer also V-Power Gasoline rated at 100 RON. Recently Agip introduced "Blu Super+", a 100 RON gasoline.
Israel: 95 RON & 98 RON are normally available at most automotive fuel stations. 96 RON is no longer available as of 2010. 95 RON is preferred because it is cheaper and performance differences are not significant. "Regular" fuel is 95 RON. All variants are unleaded.
Japan: "Regular" unleaded fuel is 90 RON and "High-octane" ("Premium") fuel is about 100 RON, or in fact 99.5 RON according to some suppliers, at least until around 2021. The minimum values are defined in standard JIS K 2202: "Regular" is >=89.0 RON, and "High-octane" is >=96.0 RON, since the revision of 1986. It means "High-octane" has a higher octane rating than the JIS standard. Although 99.5 RON is not defined, there is no significant difference in "High-octane" from different suppliers according to the president of the Petroleum Association of Japan, and it is believed that each has almost the same octane rating (99.5 RON) in spite of the JIS. But the actual octane rating is not clear and it can be sold as "High-octane" as long as it is 96.0 RON or more. "High-octane" was formerly sometimes advertised as "Octane 100", but this practice was abandoned as its actual octane value was less than 100 RON.
Latvia: 95 RON and 98 RON are widely available.
Lebanon: 95 RON and 98 RON are widely available.
Lithuania: 95 RON and 98 RON are widely available. In some gas stations E85 (bioethanol) gasoline, 98E15 (15% of ethanol), 98E25 (25% of ethanol) are available.
Malaysia: 95 RON, 97 RON and 100 RON. "Regular" unleaded fuel is 95 RON; "Premium" fuel is rated at 97 RON (Shell's V-Power Racing is rated minimum 97 RON). Petron sells 100 RON in selected outlets.
Mexico: The standard octane index is 87 AKI for regular fuel and anywhere from 91 to 93 AKI for premium fuel, although 91 AKI is the most common octane number for premium fuel. Valero is the only station offering 93 AKI fuel in Mexico, at a premium of 5% to 10% over standard 91 AKI fuel. Valero stations are usually present in main cities, such as Monterrey, Guadalajara, Querétaro and Puebla. From 1938 to 2018, Mexican government held a monopoly in the distribution of fuel, and its brands for unleaded fuel were "Pemex Magna" and "Pemex Premium", appearing in the early 1990s, before that, fuel was usually leaded. Mexican regulations do not enforce any particular labels to identify different grades of fuel as long as each grade is clearly labeled with distinct names and colors, but the long history of Pemex's colors has established a tradition of labeling regular fuel with green, premium fuel with red, and diesel with black. Gas station brands that use different colors include Shell, BP, Mobil and Akron.
Mongolia: 92 RON and 95 RON (advertised as A92 and A95 respectively) are available at nearly all stations while slightly fewer stations offer 80 RON (advertised as A80). 98 RON (advertised as A98) is available in select few stations.
Montenegro: 95 RON is sold as a "regular" fuel. As a "premium" fuel, 98 RON is sold. Both variants are unleaded.
Myanmar: Most petrol stations carry 92 RON as standard especially in rural areas. Most larger cities and highway stations have introduced 95 RON in the past few years. The highest grade available is 97 RON which is only sold by a few stations in Yangon and Nay Pyi Taw (e.g., PTT, MMTM, Petrotrans).
Netherlands: 95 RON "Euro" is sold at every station, whereas 98 RON "Super Plus" is being phased out in favor of "premium" fuels, which are all 95 RON fuels with extra additives. Shell V-Power is a 97 RON (labelled as 95 due to the legalities of only using 95 or 98 labelling), some independent tests have shown that one year after introduction it was downgraded to 95 RON, whereas in neighboring Germany Shell V-Power consists of the regular 100 RON fuel.
New Zealand: 91 RON "Regular" and 95 RON "Premium" are both widely available. 98 RON is available instead of 95 RON at some (BP, Mobil, Gull) service stations in larger urban areas (newer BP stations also offer 95 by blending 91 and 98 where 98 is available). 100 RON is available at selected NPD service stations in the South Island and in very limited locations in the North Island.
Norway: 95 RON are widely available, but 98 RON is also available at Shell, here it goes under the name v-power; it is 10-20% more expensive as 95 RON fuel. In 2023 95 RON fuel got change to 95E10 and 98 RON to 98E5. SHELL still has 98 octane under the name v-power, but also most of the Esso gas stations have 98 octane fuel as well.
Oman: 91 RON, 95 RON and 98 RON. "Regular" unleaded fuel is 91 RON; "Premium" fuel is rated at 95 RON; 98 RON in selected outlets.
Pakistan: 3 types of fuel available. 92, HOBC 95 & HOBC 97 RON. Super marketed as 92 RON, 95 RON marketed by Shell as V-Power and 97 RON by Total Parco Pakistan & Pakistan State Oil (PSO). HOBC pricing was deregulated in October, 2016.
Philippines: All automotive fuels are unleaded since December 23, 2000. Since late 2013, three grades of gasoline are available: Premium Plus, Premium (mid-grade) and Regular. Law requires the Premium Plus grade to be 97 RON or higher; Premium at 95 RON; Regular at 91 RON. Premium Plus grade fuels are exempted from having an ethanol blend, although the only Premium Plus grade available without ethanol is Petron Blaze and is rated at 100 RON. Other Premium Plus grades like Seaoil Extreme 97, Shell V-Power Racing and Unioil Gas 97 are rated at 97 RON, while Phoenix Premium 98 is rated at 98 RON. Premium grades such as Caltex Gold, Petron XCS, Phoenix Premium 95, Seaoil Extreme 95, Shell V-Power Nitro+, Unioil Gas 95 and Total Excellium are rated at 95 RON. Regular grades such as Caltex Silver, Petron Xtra Unleaded, Phoenix Super Regular 91, Seaoil Extreme U+, Shell FuelSave Unleaded, Unioil Gas 91 and Total Premier are rated at 91 RON.
Poland: Eurosuper 95 (RON 95) is sold in every gas station. Super Plus 98 (RON 98) is available in most stations, sometimes under brand (Orlen – Verva, BP – Ultimate, Shell – V-Power) and usually containing additives. Shell offers V-Power Racing fuel which is rated RON 100.
Portugal: 95 RON "Euro" is sold in every station and 98 RON "Super" being offered in almost every station.
Russia: In the Soviet Union there were different grades of automobile gasoline, which had the following names: A-56, A-66, A-70, A-72, A-74, A-76, AI-93, AI-95 also known as "Extra", and B-70 (aviation gasoline). The first letter indicated the vehicle for which the gasoline was intended, the number indicated the octane. Gasolines A-56 and A-66, A-70, and later A-72, were intended for cars with flat-head engines produced in the 1930s-1960s. Gasolines A-74, later A-76 and AI-93 for cars with overhead valve engines produced in the 1960s-1980s. AI-95 gasoline was mainly for foreign cars or government limousines ZiL and Chaika. The letter "I" in the AI-93 and AI-95 brands indicated that the octane number was calculated using the research method. After the dissolution of the Soviet Union in the 1990s, A-76 gasoline was replaced by AI-80, and AI-93 by AI-92. By the early 1980s, production of A-66 gasoline ceased, and about a decade later, so did A-72. Nowadays 92 RON is the minimum available, the standard is 95 RON is sold in every gas station. 98 RON is available in most stations. As a "premium" fuel, 100 RON is sold, Gazpromneft and Lukoil both variants are unleaded.
Saudi Arabia: Two types of fuel are available at all petrol stations in Saudi Arabia. "Premium 91" (RON 91) has green pumps, and "Super Premium 95" (RON 95) where the pumps are red. Fuel dyes are used to make the colour of the fuel match that of the pump. While petrol stations in Saudi Arabia are privatised, the prices are regulated by the authorities and have a fixed at SR 1.44 (US$0.38) and SR 2.10 (US$0.56) (as of 14 April 2019) per litre respectively; and is currently being increased at a quarterly rate to bring it up to the worldwide average by 2020. Prior to 2006, only Super Premium RON 95 was available and the pumps were not systematically coloured. The public did not know what octane rating was, so education campaigns were started, advising people to use "red petrol" only for high end cars, and to save money by using "green petrol" for regular cars and trucks.
Singapore: All four providers, Caltex, ExxonMobil, SPC and Shell have 3 grades of gasoline. Typically, these are 92, 95, and 98 RON. However, since 2009, Shell has removed 92 RON.
South Africa: "regular" unleaded fuel is 95 RON in coastal areas. Inland (higher elevation) "regular" unleaded fuel is 93 RON; once again most fuel stations optionally offer 95 RON.
South Korea: "regular" unleaded fuel is 91~94 RON, "premium" is 95+ RON nationally. However, not all gas stations carry "premium."
Spain: 95 RON "Euro" is sold in every station with 98 RON "Super" being offered in most stations. Many stations around cities and highways offer other high-octane "premium" brands.
Sri Lanka: Sri Lanka switched their regular gasoline from 90 RON to 92 RON on January 1, 2014. In Ceypetco filling stations, 92 RON is the regular automotive fuel and 95 RON is called 'Super Petrol', which comes at a premium price. In LIOC filling stations, 92 RON is the regular automotive fuel and 95 RON is available as 'Premium Petrol'. As of 2022, LIOC fillings stations offer a new fuel labelled as 'XtraPremium' Petrol which is marketed as 'Euro 3' standard petrol. Similarly 95 RON petrol is offered as 'XtraPremium' 95 Petrol. Sri Lanka adopted RON 100 Octane 100 from July 2024. It is the 8th country in the world to use RON 100.
Sweden: 95 RON, 98 RON and E85 are widely available.
Taiwan: 92 RON, 95 RON and 98 RON are widely available at gas stations in Taiwan.
Thailand: 95 RON E0, 95 RON E10, 91 RON E10, 95 RON E20 are widely available in all parts of Thailand. 97 RON E10 fuel is also available in some Bangchak's filling stations in many parts of Thailand.
Trinidad and Tobago: 92 RON (Super) and 95 RON (Premium) are widely available.
Turkey: 95 RON and 95+ RON widely available in gas stations. 91 RON (Regular) has been dropped in 2006. 98 and 100 RON (Shell V-Power Racing) has been dropped in late 2009. The Gas which has been advertised 97 RON has been dropped in 2014 and renamed 95+.
Ukraine: 80 RON and 98 RON gasoline is available. The standard gasoline is 95 RON, but 92 RON gasoline is also widely available and popular for older cars. There is no government regulation for gasoline with RON higher than 98 so some stations are marketing 100 RON gasoline when in reality this can be anything above 98 RON with extra cleaning additives.
United Kingdom: 'regular' gasoline has an octane rating of 95 RON, with 97 RON fuel being widely available as the Super Unleaded. Tesco and Shell both offer 99 RON fuel. In April 2006, BP started a public trial of the super-high octane gasoline BP Ultimate Unleaded 102, which as the name suggests, has an octane rating of 102 RON. Although BP Ultimate Unleaded (with an octane rating of 97 RON) and BP Ultimate Diesel are both widely available throughout the UK, BP Ultimate Unleaded 102 was available throughout the UK in only 10 filling stations, and was priced at about two and half times more than their 97 RON fuel. In March 2010, BP stopped sales of Ultimate Unleaded 102, citing the closure of their specialty fuels manufacturing facility. Shell V-Power is also available, but in a 99 RON octane rating, and Tesco fuel stations also supply the Greenergy produced 99 RON "Momentum99".
United States: in the US octane rating is displayed in AKI. In most areas, the standard grades are 87, 89–90, and 91–94 AKI. In the Rocky Mountain (high elevation) states, 85 AKI (90 RON) is the minimum octane, and 91 AKI (95 RON) is the maximum octane available in fuel. The reason for this is that in higher-elevation areas, a typical naturally aspirated engine draws in less air mass per cycle because of the reduced density of the atmosphere. This directly translates to less fuel and reduced absolute compression in the cylinder, therefore deterring knock. It is safe to fill a carbureted car that normally takes 87 AKI fuel at sea level with 85 AKI fuel in the mountains, but at sea level the fuel may cause damage to the engine. Fuel injectors have almost completely replaced carburettors in nearly all modern automobiles produced from the late 1980s and early 1990s onwards. 85 AKI fuel is not recommended for modern automobiles and may cause damage to the engine and decreased performance. Another disadvantage to this strategy is that most turbocharged vehicles are unable to produce full power, even when using the "premium" 91 AKI fuel. In some east coast states, up to 94 AKI (99 or 100 RON) is available. As of January 2011, over 40 states and a total of over 2500 stations offer ethanol-based E-85 fuel with 94–96 AKI. Often, filling stations near US racing tracks will offer higher octane levels such as 100 AKI from separate dedicated pumps.
State standard gasoline grades: U.S. State Fuel Octane Standards
Venezuela: 91 RON and 95 RON gasoline is available nationwide, in all PDV gas stations. 95 RON gasoline is the most widely used in the country, although most cars in Venezuela would work with 91 RON gasoline. This is because gasoline prices are heavily subsidized by the government (US$0.0083 per gallon 95 RON, vs US$0.0061 per gallon 91 RON). All gasoline in Venezuela is unleaded.
Vietnam: 92 RON is in every gas station and 95 RON is in the urban area. They start selling A92-E5 gasoline (92 RON with 5 percent of ethanol) at 2017. On January 1, 2018, Vietnamese government forced every gas station stop selling 92 RON and sell 95 RON + A92-E5 instead. From 2022, Vietnam will start selling gasoline according to Euro 5 standards, with the choices 95 RON and 97 RON(in SFC gas stations).
Zimbabwe: 93 octane available with no other grades of fuels available, E10 which is an ethanol blend of fuel at 10% ethanol is available the octane rating however is still to be tested and confirmed but it is assumed that it is around 95 Octane. E85 available from 3 outlets with an octane rating AKI index of between 102 and 105 depending on the base gasoline the ethanol is blended with.
Misconceptions around octane rating
Due to its name, the chemical "octane" is often misunderstood as the only substance that determines the octane rating (or octane number) of a fuel. This is an inaccurate description. In reality, the octane rating is defined as a number describing the stability and ability of a fuel to prevent an engine from unwanted combustions that occur spontaneously in the other regions within a cylinder (i.e., delocalized explosions from the spark plug). This phenomenon of combustion is more commonly known as engine knocking or self-ignition, which causes damage to pistons over time and reduces the lifespan of engines.
In 1927, Graham Edgar devised the method of using iso-octane and n-heptane as reference chemicals, in order to rate the knock resistance of a fuel with respect to this isomer of octane, thus the name "octane rating". By definition, the isomers iso-octane and n-heptane have an octane rating of 100 and 0, respectively. Because of its more volatile nature, n-heptane ignites and knocks readily, which gives it a relatively low octane rating; the isomer iso-octane causes less knocking because it is more branched and combusts more smoothly. In general, branched compounds with a higher intermolecular force (e.g., London dispersion force for iso-octane) will have a higher octane rating, because they are harder to ignite.
Octane ratings of octane isomers
Octane isomers such as n-octane and 2,3,3-trimethylpentane have an octane rating of -20 and 106.1, respectively (RON measurement). The large differences between the octane ratings for the isomers show that the compound octane itself is clearly not the only factor that determines octane ratings, especially for commercial fuels consist of a wide variety of compounds.
Octane in culture
"Octane" is colloquially used in the expression "high-octane". The term is used to describe a powerful action because of the association with the concept of "octane rating". This is a misleading term, because the octane rating of gasoline is not directly related to the power output of an engine. Using gasoline of a higher octane than an engine is designed for cannot increase its power output.
Octane became well known in American popular culture in the 1960s, when gasoline companies boasted of "high octane" levels in their gasoline advertisements. The compound adjective "high-octane", meaning powerful or dynamic, is recorded in a figurative sense from 1944. By the 1990s, the phrase was commonly being used as a word intensifier, and it has found a place in modern English slang.
| Physical sciences | Hydrocarbons | Chemistry |
147252 | https://en.wikipedia.org/wiki/Integration%20by%20parts | Integration by parts | In calculus, and more generally in mathematical analysis, integration by parts or partial integration is a process that finds the integral of a product of functions in terms of the integral of the product of their derivative and antiderivative. It is frequently used to transform the antiderivative of a product of functions into an antiderivative for which a solution can be more easily found. The rule can be thought of as an integral version of the product rule of differentiation; it is indeed derived using the product rule.
The integration by parts formula states:
Or, letting and while and the formula can be written more compactly:
The former expression is written as a definite integral and the latter is written as an indefinite integral. Applying the appropriate limits to the latter expression should yield the former, but the latter is not necessarily equivalent to the former.
Mathematician Brook Taylor discovered integration by parts, first publishing the idea in 1715. More general formulations of integration by parts exist for the Riemann–Stieltjes and Lebesgue–Stieltjes integrals. The discrete analogue for sequences is called summation by parts.
Theorem
Product of two functions
The theorem can be derived as follows. For two continuously differentiable functions and , the product rule states:
Integrating both sides with respect to ,
and noting that an indefinite integral is an antiderivative gives
where we neglect writing the constant of integration. This yields the formula for integration by parts:
or in terms of the differentials ,
This is to be understood as an equality of functions with an unspecified constant added to each side. Taking the difference of each side between two values and and applying the fundamental theorem of calculus gives the definite integral version:
The original integral contains the derivative ; to apply the theorem, one must find , the antiderivative of , then evaluate the resulting integral
Validity for less smooth functions
It is not necessary for and to be continuously differentiable. Integration by parts works if is absolutely continuous and the function designated is Lebesgue integrable (but not necessarily continuous). (If has a point of discontinuity then its antiderivative may not have a derivative at that point.)
If the interval of integration is not compact, then it is not necessary for to be absolutely continuous in the whole interval or for to be Lebesgue integrable in the interval, as a couple of examples (in which and are continuous and continuously differentiable) will show. For instance, if
is not absolutely continuous on the interval , but nevertheless:
so long as is taken to mean the limit of as and so long as the two terms on the right-hand side are finite. This is only true if we choose Similarly, if
is not Lebesgue integrable on the interval , but nevertheless
with the same interpretation.
One can also easily come up with similar examples in which and are not continuously differentiable.
Further, if is a function of bounded variation on the segment and is differentiable on then
where denotes the signed measure corresponding to the function of bounded variation , and functions are extensions of to which are respectively of bounded variation and differentiable.
Product of many functions
Integrating the product rule for three multiplied functions, , , , gives a similar result:
In general, for factors
which leads to
Visualization
Consider a parametric curve . Assuming that the curve is locally one-to-one and integrable, we can define
The area of the blue region is
Similarly, the area of the red region is
The total area A1 + A2 is equal to the area of the bigger rectangle, x2y2, minus the area of the smaller one, x1y1:
Or, in terms of t,
Or, in terms of indefinite integrals, this can be written as
Rearranging:
Thus integration by parts may be thought of as deriving the area of the blue region from the area of rectangles and that of the red region.
This visualization also explains why integration by parts may help find the integral of an inverse function f−1(x) when the integral of the function f(x) is known. Indeed, the functions x(y) and y(x) are inverses, and the integral ∫ x dy may be calculated as above from knowing the integral ∫ y dx. In particular, this explains use of integration by parts to integrate logarithm and inverse trigonometric functions. In fact, if is a differentiable one-to-one function on an interval, then integration by parts can be used to derive a formula for the integral of in terms of the integral of . This is demonstrated in the article, Integral of inverse functions.
Applications
Finding antiderivatives
Integration by parts is a heuristic rather than a purely mechanical process for solving integrals; given a single function to integrate, the typical strategy is to carefully separate this single function into a product of two functions u(x)v(x) such that the residual integral from the integration by parts formula is easier to evaluate than the single function. The following form is useful in illustrating the best strategy to take:
On the right-hand side, u is differentiated and v is integrated; consequently it is useful to choose u as a function that simplifies when differentiated, or to choose v as a function that simplifies when integrated. As a simple example, consider:
Since the derivative of ln(x) is , one makes (ln(x)) part u; since the antiderivative of is −, one makes part v. The formula now yields:
The antiderivative of − can be found with the power rule and is .
Alternatively, one may choose u and v such that the product u′ (∫v dx) simplifies due to cancellation. For example, suppose one wishes to integrate:
If we choose u(x) = ln(|sin(x)|) and v(x) = sec2x, then u differentiates to using the chain rule and v integrates to tan x; so the formula gives:
The integrand simplifies to 1, so the antiderivative is x. Finding a simplifying combination frequently involves experimentation.
In some applications, it may not be necessary to ensure that the integral produced by integration by parts has a simple form; for example, in numerical analysis, it may suffice that it has small magnitude and so contributes only a small error term. Some other special techniques are demonstrated in the examples below.
Polynomials and trigonometric functions
In order to calculate
let:
then:
where C is a constant of integration.
For higher powers of in the form
repeatedly using integration by parts can evaluate integrals such as these; each application of the theorem lowers the power of by one.
Exponentials and trigonometric functions
An example commonly used to examine the workings of integration by parts is
Here, integration by parts is performed twice. First let
then:
Now, to evaluate the remaining integral, we use integration by parts again, with:
Then:
Putting these together,
The same integral shows up on both sides of this equation. The integral can simply be added to both sides to get
which rearranges to
where again (and ) is a constant of integration.
A similar method is used to find the integral of secant cubed.
Functions multiplied by unity
Two other well-known examples are when integration by parts is applied to a function expressed as a product of 1 and itself. This works if the derivative of the function is known, and the integral of this derivative times is also known.
The first example is . We write this as:
Let:
then:
where is the constant of integration.
The second example is the inverse tangent function :
Rewrite this as
Now let:
then
using a combination of the inverse chain rule method and the natural logarithm integral condition.
LIATE rule
The LIATE rule is a rule of thumb for integration by parts. It involves choosing as u the function that comes first in the following list:
L – logarithmic functions: etc.
I – inverse trigonometric functions (including hyperbolic analogues): etc.
A – algebraic functions (such as polynomials): etc.
T – trigonometric functions (including hyperbolic analogues): etc.
E – exponential functions: etc.
The function which is to be dv is whichever comes last in the list. The reason is that functions lower on the list generally have simpler antiderivatives than the functions above them. The rule is sometimes written as "DETAIL", where D stands for dv and the top of the list is the function chosen to be dv. An alternative to this rule is the ILATE rule, where inverse trigonometric functions come before logarithmic functions.
To demonstrate the LIATE rule, consider the integral
Following the LIATE rule, u = x, and dv = cos(x) dx, hence du = dx, and v = sin(x), which makes the integral become
which equals
In general, one tries to choose u and dv such that du is simpler than u and dv is easy to integrate. If instead cos(x) was chosen as u, and x dx as dv, we would have the integral
which, after recursive application of the integration by parts formula, would clearly result in an infinite recursion and lead nowhere.
Although a useful rule of thumb, there are exceptions to the LIATE rule. A common alternative is to consider the rules in the "ILATE" order instead. Also, in some cases, polynomial terms need to be split in non-trivial ways. For example, to integrate
one would set
so that
Then
Finally, this results in
Integration by parts is often used as a tool to prove theorems in mathematical analysis.
Wallis product
The Wallis infinite product for
may be derived using integration by parts.
Gamma function identity
The gamma function is an example of a special function, defined as an improper integral for . Integration by parts illustrates it to be an extension of the factorial function:
Since
when is a natural number, that is, , applying this formula repeatedly gives the factorial:
Use in harmonic analysis
Integration by parts is often used in harmonic analysis, particularly Fourier analysis, to show that quickly oscillating integrals with sufficiently smooth integrands decay quickly. The most common example of this is its use in showing that the decay of function's Fourier transform depends on the smoothness of that function, as described below.
Fourier transform of derivative
If is a -times continuously differentiable function and all derivatives up to the th one decay to zero at infinity, then its Fourier transform satisfies
where is the th derivative of . (The exact constant on the right depends on the convention of the Fourier transform used.) This is proved by noting that
so using integration by parts on the Fourier transform of the derivative we get
Applying this inductively gives the result for general . A similar method can be used to find the Laplace transform of a derivative of a function.
Decay of Fourier transform
The above result tells us about the decay of the Fourier transform, since it follows that if and are integrable then
In other words, if satisfies these conditions then its Fourier transform decays at infinity at least as quickly as . In particular, if then the Fourier transform is integrable.
The proof uses the fact, which is immediate from the definition of the Fourier transform, that
Using the same idea on the equality stated at the start of this subsection gives
Summing these two inequalities and then dividing by gives the stated inequality.
Use in operator theory
One use of integration by parts in operator theory is that it shows that the (where ∆ is the Laplace operator) is a positive operator on (see Lp space). If is smooth and compactly supported then, using integration by parts, we have
Other applications
Determining boundary conditions in Sturm–Liouville theory
Deriving the Euler–Lagrange equation in the calculus of variations
Repeated integration by parts
Considering a second derivative of in the integral on the LHS of the formula for partial integration suggests a repeated application to the integral on the RHS:
Extending this concept of repeated partial integration to derivatives of degree leads to
This concept may be useful when the successive integrals of are readily available (e.g., plain exponentials or sine and cosine, as in Laplace or Fourier transforms), and when the th derivative of vanishes (e.g., as a polynomial function with degree ). The latter condition stops the repeating of partial integration, because the RHS-integral vanishes.
In the course of the above repetition of partial integrations the integrals
and and
get related. This may be interpreted as arbitrarily "shifting" derivatives between and within the integrand, and proves useful, too (see Rodrigues' formula).
Tabular integration by parts
The essential process of the above formula can be summarized in a table; the resulting method is called "tabular integration" and was featured in the film Stand and Deliver (1988).
For example, consider the integral
and take
Begin to list in column A the function and its subsequent derivatives until zero is reached. Then list in column B the function and its subsequent integrals until the size of column B is the same as that of column A. The result is as follows:
{| class="wikitable" style="text-align:center"
!# i !! Sign !! A: derivatives !! B: integrals
|-
| 0 || + || ||
|-
| 1 || − || ||
|-
| 2 || + || ||
|-
| 3 || − || ||
|-
| 4 || + || ||
|}
The product of the entries in of columns A and B together with the respective sign give the relevant integrals in in the course of repeated integration by parts. yields the original integral. For the complete result in the must be added to all the previous products () of the of column A and the of column B (i.e., multiply the 1st entry of column A with the 2nd entry of column B, the 2nd entry of column A with the 3rd entry of column B, etc. ...) with the given This process comes to a natural halt, when the product, which yields the integral, is zero ( in the example). The complete result is the following (with the alternating signs in each term):
This yields
The repeated partial integration also turns out useful, when in the course of respectively differentiating and integrating the functions and their product results in a multiple of the original integrand. In this case the repetition may also be terminated with this index This can happen, expectably, with exponentials and trigonometric functions. As an example consider
{| class="wikitable" style="text-align:center"
!# i !! Sign !! A: derivatives !! B: integrals
|-
| 0 || + || ||
|-
| 1 || − || ||
|-
| 2 || + || ||
|}
In this case the product of the terms in columns A and B with the appropriate sign for index yields the negative of the original integrand (compare
Observing that the integral on the RHS can have its own constant of integration , and bringing the abstract integral to the other side, gives:
and finally:
where .
Higher dimensions
Integration by parts can be extended to functions of several variables by applying a version of the fundamental theorem of calculus to an appropriate product rule. There are several such pairings possible in multivariate calculus, involving a scalar-valued function u and vector-valued function (vector field) V.
The product rule for divergence states:
Suppose is an open bounded subset of with a piecewise smooth boundary . Integrating over with respect to the standard volume form , and applying the divergence theorem, gives:
where is the outward unit normal vector to the boundary, integrated with respect to its standard Riemannian volume form . Rearranging gives:
or in other words
The regularity requirements of the theorem can be relaxed. For instance, the boundary need only be Lipschitz continuous, and the functions u, v need only lie in the Sobolev space .
Green's first identity
Consider the continuously differentiable vector fields and , where is the i-th standard basis vector for . Now apply the above integration by parts to each times the vector field :
Summing over i gives a new integration by parts formula:
The case , where , is known as the first of Green's identities:
| Mathematics | Integral calculus | null |
147341 | https://en.wikipedia.org/wiki/Bog | Bog | A bog or bogland is a wetland that accumulates peat as a deposit of dead plant materials often mosses, typically sphagnum moss. It is one of the four main types of wetlands. Other names for bogs include mire, mosses, quagmire, and muskeg; alkaline mires are called fens. A bayhead is another type of bog found in the forest of the Gulf Coast states in the United States. They are often covered in heath or heather shrubs rooted in the sphagnum moss and peat. The gradual accumulation of decayed plant material in a bog functions as a carbon sink.
Bogs occur where the water at the ground surface is acidic and low in nutrients. A bog usually is found at a freshwater soft spongy ground that is made up of decayed plant matter which is known as peat. They are generally found in cooler northern climates and are formed in poorly draining lake basins. In contrast to fens, they derive most of their water from precipitation rather than mineral-rich ground or surface water. Water flowing out of bogs has a characteristic brown colour, which comes from dissolved peat tannins. In general, the low fertility and cool climate result in relatively slow plant growth, but decay is even slower due to low oxygen levels in saturated bog soils. Hence, peat accumulates. Large areas of the landscape can be covered many meters deep in peat.
Bogs have distinctive assemblages of animal, fungal, and plant species, and are of high importance for biodiversity, particularly in landscapes that are otherwise settled and farmed.
Distribution and extent
Bogs are widely distributed in cold, temperate climes, mostly in boreal ecosystems in the Northern Hemisphere. The world's largest wetland is the peat bogs of the Western Siberian Lowlands in Russia, which cover more than a million square kilometres. Large peat bogs also occur in North America, particularly the Hudson Bay Lowland and the Mackenzie River Basin. They are less common in the Southern Hemisphere, with the largest being the Magellanic moorland, comprising some in southern South America. Sphagnum bogs were widespread in northern Europe but have often been cleared and drained for agriculture. A paper led by Graeme T. Swindles in 2019 showed that peatlands across Europe have undergone rapid drying in recent centuries owing to human impacts including drainage, peat cutting and burning.
A 2014 expedition leaving from Itanga village, Republic of the Congo, discovered a peat bog "as big as England" which stretches into neighboring Democratic Republic of Congo.
Definition
Like all wetlands, it is difficult to rigidly define bogs for a number of reasons, including variations between bogs, the in-between nature of wetlands as an intermediate between terrestrial and aquatic ecosystems, and varying definitions between wetland classification systems. However, there are characteristics common to all bogs that provide a broad definition:
Peat is present, usually thicker than .
The wetland receives most of its water and nutrients from precipitation (ombrotrophic) rather than surface or groundwater (minerotrophic).
The wetland is nutrient-poor (oligotrophic).
The wetland is strongly acidic (bogs near coastal areas may be less acidic due to sea spray).
Because all bogs have peat, they are a type of peatland. As a peat-producing ecosystem, they are also classified as mires, along with fens. Bogs differ from fens, in that fens receive water and nutrients from mineral-rich surface or groundwater, while bogs receive water and nutrients from precipitation. Because fens are supplied with mineral-rich water, they tend to range from slightly acidic to slightly basic, while bogs are always acidic because precipitation lacks the dissolved minerals (e.g. calcium, magnesium, carbonate) that act to buffer the natural acidity of atmospheric carbon dioxide. Geography and geology both impact the hydrology: as groundwater mineral content reflects the bedrock geology, there can be great deal of variability in some common ions (e.g. manganese, iron) while proximity to coastal areas is associated with higher sulfate and sodium concentrations.
Ecology and protection
There are many highly specialized animals, fungi, and plants associated with bog habitat. Most are capable of tolerating the combination of low nutrient levels and waterlogging. Sphagnum is generally abundant, along with ericaceous shrubs. The shrubs are often evergreen, which may assist in conservation of nutrients. In drier locations, evergreen trees can occur, in which case the bog blends into the surrounding expanses of boreal evergreen forest. Sedges are one of the more common herbaceous species. Carnivorous plants such as sundews (Drosera) and pitcher plants (for example Sarracenia purpurea) have adapted to the low-nutrient conditions by using invertebrates as a nutrient source. Orchids have adapted to these conditions through the use of mycorrhizal fungi to extract nutrients. Some shrubs such as Myrica gale (bog myrtle) have root nodules in which nitrogen fixation occurs, thereby providing another supplemental source of nitrogen.
Bogs are recognized as a significant/specific habitat type by a number of governmental and conservation agencies. They can provide habitat for mammals, such as caribou, moose, and beavers, as well as for species of nesting shorebirds, such as Siberian cranes and yellowlegs. Bogs contain species of vulnerable reptilians such as the bog turtle. Bogs even have distinctive insects; English bogs give a home to a yellow fly called the hairy canary fly (Phaonia jaroschewskii), and bogs in North America are habitat for a butterfly called the bog copper (Lycaena epixanthe). In Ireland, the viviparous lizard, the only known reptile in the country, dwells in bogland.
The United Kingdom in its Biodiversity Action Plan establishes bog habitats as a priority for conservation. Russia has a large reserve system in the West Siberian Lowland. The highest protected status occurs in Zapovedniks (IUCN category IV); Gydansky and Yugansky are two prominent examples.
Bogs are fragile ecosystems, and have been deteriorating quickly, as archaeologists and scientists have been recently finding. Bone material found in bogs has had accelerated deterioration from first analyses in the 1940s. This has been found to be from fluctuations in ground water and increase in acidity in lower areas of bogs that is affecting the rich organic material. Many of these areas have been permeated to the lowest levels with oxygen, which dries and cracks layers. There have been some temporary solutions to try and fix these issues, such as adding soil to the tops of threatened areas, yet they do not work in the long-term. Extreme weather like dry summers are likely the cause, as they lower precipitation and the groundwater table. It is speculated that these issues will only increase with a rise in global temperature and climate change. Since bogs take thousands of years to form and create the rich peat that is used as a resource, once they are gone they are extremely hard to recover. Arctic and sub-Arctic circles where many bogs are warming at 0.6 °C per decade, an amount twice as large as the global average. Because bogs and other peatlands are carbon sinks, they are releasing large amounts of greenhouse gases as they warm up. These changes have resulted in a severe decline of biodiversity and species populations of peatlands throughout Northern Europe.
Types
Bog habitats may develop in various situations, depending on the climate and topography.
By location and water source
Bogs may be classified on their topography, proximity to water, method of recharge, and nutrient accumulation.
Valley bog
These develop in gently sloping valleys or hollows. A layer of peat fills the deepest part of the valley, and a stream may run through the surface of the bog. Valley bogs may develop in relatively dry and warm climates, but because they rely on ground or surface water, they only occur on acidic substrates.
Raised bog
These develop from a lake or flat marshy area, over either non-acidic or acidic substrates. Over centuries there is a progression from open lake, to a marsh, to a fen (or, on acidic substrates, valley bog), to a carr, as silt or peat accumulates within the lake. Eventually, peat builds up to a level where the land surface is too flat for ground or surface water to reach the center of the wetland. This part, therefore, becomes wholly rain-fed (ombrotrophic), and the resulting acidic conditions allow the development of bog (even if the substrate is non-acidic). The bog continues to form peat, and over time a shallow dome of bog peat develops into a raised bog. The dome is typically a few meters high in the center and is often surrounded by strips of fen or other wetland vegetation at the edges or along streamsides where groundwater can percolate into the wetland.
The various types of raised bog may be divided into:
Coastal bog
Plateau bog
Upland bog
Kermi bog
String bog
Palsa bog
Polygonal bog
Blanket bog
In cool climates with consistently high rainfall (on more than c. 235 days a year), the ground surface may remain waterlogged for much of the time, providing conditions for the development of bog vegetation. In these circumstances, bog develops as a layer "blanketing" much of the land, including hilltops and slopes. Although a blanket bog is more common on acidic substrates, under some conditions it may also develop on neutral or even alkaline ones, if abundant acidic rainwater predominates over the groundwater. A blanket bog can occur in drier or warmer climates, because under those conditions hilltops and sloping ground dry out too often for peat to form – in intermediate climates a blanket bog may be limited to areas which are shaded from direct sunshine. In periglacial climates a patterned form of blanket bog may occur, known as a string bog. In Europe, these mostly very thin peat layers without significant surface structures are distributed over the hills and valleys of Ireland, Scotland, England, and Norway. In North America, blanket bogs occur predominantly in Canada east of Hudson Bay. These bogs are often still under the influence of mineral soil water (groundwater). Blanket bogs do not occur north of the 65th latitude in the northern hemisphere.
Quaking bog
A quaking bog, schwingmoor, or swingmoor is a form of floating bog occurring in wetter parts of valley bogs and raised bogs and sometimes around the edges of acidic lakes. The bog vegetation, mostly sphagnum moss anchored by sedges (such as Carex lasiocarpa), forms a floating mat approximately half a meter thick on the surface of water or above very wet peat. White spruce (Picea glauca) may grow in this bog regime. Walking on the surface causes it to move – larger movements may cause visible ripples on the surface, or they may even make trees sway. The bog mat may eventually spread across the water surface to cover bays or even entire small lakes. Bogs at the edges of lakes may become detached and form floating islands.
Cataract bog
A cataract bog is a rare ecological community formed where a permanent stream flows over a granite outcropping. The sheeting of water keeps the edges of the rock wet without eroding the soil, but in this precarious location, no tree or large shrub can maintain a roothold. The result is a narrow, permanently wet habitat.
Uses
Industrial uses
After drying, peat is used as a fuel, and it has been used that way for centuries. More than 20% of home heat in Ireland comes from peat, and it is also used for fuel in Finland, Scotland, Germany, and Russia. Russia is the leading exporter of peat for fuel, at more than 90 million metric tons per year. Ireland's Bord na Móna ("peat board") was one of the first companies to mechanically harvest peat, which is being phased out.
The other major use of dried peat is as a soil amendment (sold as moss peat or sphagnum peat) to increase the soil's capacity to retain moisture and enrich the soil. It is also used as a mulch. Some distilleries, notably in the Islay whisky-producing region, use the smoke from peat fires to dry the barley used in making Scotch whisky.
Once the peat has been extracted it can be difficult to restore the wetland, since peat accumulation is a slow process. More than 90% of the bogs in England have been damaged or destroyed. In 2011 plans for the elimination of peat in gardening products were announced by the UK government.
Other uses
The peat in bogs is an important place for the storage of carbon. If the peat decays, carbon dioxide would be released to the atmosphere, contributing to global warming. Undisturbed, bogs function as a carbon sink. As one example, the peatlands of the former Soviet Union were calculated to be removing 52 Tg of carbon per year from the atmosphere. Therefore, the rewetting of drained peatlands may be one of the most cost-effective ways to mitigate climate change.
Peat bogs are also important in storing fresh water, particularly in the headwaters of large rivers. Even the enormous Yangtze River arises in the Ruoergai peatland near its headwaters in Tibet.
Blueberries, cranberries, cloudberries, huckleberries, and lingonberries are harvested from the wild in bogs. Bog oak, wood that has been partially preserved by bogs, has been used in the manufacture of furniture.
Sphagnum bogs are also used for outdoor recreation, with activities including ecotourism and hunting. For example, many popular canoe routes in northern Canada include areas of peatland. Some other activities, such as all-terrain vehicle use, are especially damaging to bogs.
Archaeology
The anaerobic environment and presence of tannic acids within bogs can result in the remarkable preservation of organic material. Finds of such material have been made in Slovenia, Denmark, Germany, Ireland, Russia, and the United Kingdom. Some bogs have preserved bog-wood, such as ancient oak logs useful in dendrochronology. They have yielded extremely well-preserved bog bodies, with hair, organs, and skin intact, buried there thousands of years ago after apparent Germanic and Celtic human sacrifice. Excellent examples of such human specimens include the Haraldskær Woman and Tollund Man in Denmark, and Lindow man found at Lindow Common in England. The Tollund Man was so well preserved that when the body was discovered in 1950, the discoverers thought it was a recent murder victim and researchers were even able to tell the last meal that the Tollund Man ate before he died: porridge and fish. This process happens because of the low oxygen levels of bogs in combination with the high acidity. These anaerobic conditions lead to some of the best-preserved mummies and offer much archeological insight into society as far as 8,000 years back. Céide Fields in County Mayo in Ireland, a 5,000-year-old neolithic farming landscape has been found preserved under a blanket bog, complete with field walls and hut sites. One ancient artifact found in various bogs is bog butter, large masses of fat, usually in wooden containers. These are thought to have been food stores of both butter and tallow.
Image gallery
| Physical sciences | Wetlands | null |
147460 | https://en.wikipedia.org/wiki/Free%20variables%20and%20bound%20variables | Free variables and bound variables | In mathematics, and in other disciplines involving formal languages, including mathematical logic and computer science, a variable may be said to be either free or bound. Some older books use the terms real variable and apparent variable for free variable and bound variable, respectively. A free variable is a notation (symbol) that specifies places in an expression where substitution may take place and is not a parameter of this or any container expression. The idea is related to a placeholder (a symbol that will later be replaced by some value), or a wildcard character that stands for an unspecified symbol.
In computer programming, the term free variable refers to variables used in a function that are neither local variables nor parameters of that function. The term non-local variable is often a synonym in this context.
An instance of a variable symbol is bound, in contrast, if the value of that variable symbol has been bound to a specific value or range of values in the domain of discourse or universe. This may be achieved through the use of logical quantifiers, variable-binding operators, or an explicit statement of allowed values for the variable (such as, "...where is a positive integer".) A variable symbol overall is bound if at least one occurrence of it is bound.pp.142--143 Since the same variable symbol may appear in multiple places in an expression, some occurrences of the variable symbol may be free while others are bound,p.78 hence "free" and "bound" are at first defined for occurrences and then generalized over all occurrences of said variable symbol in the expression. However it is done, the variable ceases to be an independent variable on which the value of the expression depends, whether that value be a truth value or the numerical result of a calculation, or, more generally, an element of an image set of a function.
While the domain of discourse in many contexts is understood, when an explicit range of values for the bound variable has not been given, it may be necessary to specify the domain in order to properly evaluate the expression. For example, consider the following expression in which both variables are bound by logical quantifiers:
This expression evaluates to false if the domain of and is the real numbers, but true if the domain is the complex numbers.
The term "dummy variable" is also sometimes used for a bound variable (more commonly in general mathematics than in computer science), but this should not be confused with the identically named but unrelated concept of dummy variable as used in statistics, most commonly in regression analysis.p.17
Examples
Before stating a precise definition of free variable and bound variable, the following are some examples that perhaps make these two concepts clearer than the definition would:
In the expression
n is a free variable and k is a bound variable; consequently the value of this expression depends on the value of n, but there is nothing called k on which it could depend.
In the expression
y is a free variable and x is a bound variable; consequently the value of this expression depends on the value of y, but there is nothing called x on which it could depend.
In the expression
x is a free variable and h is a bound variable; consequently the value of this expression depends on the value of x, but there is nothing called h on which it could depend.
In the expression
z is a free variable and x and y are bound variables, associated with logical quantifiers; consequently the logical value of this expression depends on the value of z, but there is nothing called x or y on which it could depend.
More widely, in most proofs, bound variables are used. For example, the following proof shows that all squares of positive even integers are divisible by
Let be a positive even integer. Then there is an integer such that . Since , we have divisible by
not only k but also n have been used as bound variables as a whole in the proof.
Variable-binding operators
The following
are some common variable-binding operators. Each of them binds the variable x for some set S.
Many of these are operators which act on functions of the bound variable. In more complicated contexts, such notations can become awkward and confusing. It can be useful to switch to notations which make the binding explicit, such as
for sums or
for differentiation.
Formal explanation
Variable-binding mechanisms occur in different contexts in mathematics, logic and computer science. In all cases, however, they are purely syntactic properties of expressions and variables in them. For this section we can summarize syntax by identifying an expression with a tree whose leaf nodes are variables, constants, function constants or predicate constants and whose non-leaf nodes are logical operators. This expression can then be determined by doing an inorder traversal of the tree. Variable-binding operators are logical operators that occur in almost every formal language. A binding operator Q takes two arguments: a variable v and an expression P, and when applied to its arguments produces a new expression Q(v, P). The meaning of binding operators is supplied by the semantics of the language and does not concern us here.
Variable binding relates three things: a variable v, a location a for that variable in an expression and a non-leaf node n of the form Q(v, P). Note: we define a location in an expression as a leaf node in the syntax tree. Variable binding occurs when that location is below the node n.
In the lambda calculus, x is a bound variable in the term M = λx. T and a free variable in the term T. We say x is bound in M and free in T. If T contains a subterm λx. U then x is rebound in this term. This nested, inner binding of x is said to "shadow" the outer binding. Occurrences of x in U are free occurrences of the new x.
Variables bound at the top level of a program are technically free variables within the terms to which they are bound but are often treated specially because they can be compiled as fixed addresses. Similarly, an identifier bound to a recursive function is also technically a free variable within its own body but is treated specially.
A closed term is one containing no free variables.
Function expressions
To give an example from mathematics, consider an expression which defines a function
where t is an expression. t may contain some, all or none of the x1, …, xn and it may contain other variables. In this case we say that function definition binds the variables x1, …, xn.
In this manner, function definition expressions of the kind shown above can be thought of as the variable binding operator, analogous to the lambda expressions of lambda calculus. Other binding operators, like the summation sign, can be thought of as higher-order functions applying to a function. So, for example, the expression
could be treated as a notation for
where is an operator with two parameters—a one-parameter function, and a set to evaluate that function over. The other operators listed above can be expressed in similar ways; for example, the universal quantifier can be thought of as an operator that evaluates to the logical conjunction of the Boolean-valued function P applied over the (possibly infinite) set S.
Natural language
When analyzed in formal semantics, natural languages can be seen to have free and bound variables. In English, personal pronouns like he, she, they, etc. can act as free variables.
Lisa found her book.
In the sentence above, the possessive pronoun her is a free variable. It may refer to the previously mentioned Lisa or to any other female. In other words, her book could be referring to Lisa's book (an instance of coreference) or to a book that belongs to a different female (e.g. Jane's book). Whoever the referent of her is can be established according to the situational (i.e. pragmatic) context. The identity of the referent can be shown using coindexing subscripts where i indicates one referent and j indicates a second referent (different from i). Thus, the sentence Lisa found her book has the following interpretations:
Lisai found heri book. (interpretation #1: her = of Lisa)
Lisai found herj book. (interpretation #2: her = of a female that is not Lisa)
The distinction is not purely of academic interest, as some languages do actually have different forms for heri and herj: for example, Norwegian and Swedish translate coreferent heri as sin and noncoreferent herj as hennes.
English does allow specifying coreference, but it is optional, as both interpretations of the previous example are valid (the ungrammatical interpretation is indicated with an asterisk):
Lisai found heri own book. (interpretation #1: her = of Lisa)
*Lisai found herj own book. (interpretation #2: her = of a female that is not Lisa)
However, reflexive pronouns, such as himself, herself, themselves, etc., and reciprocal pronouns, such as each other, act as bound variables. In a sentence like the following:
Jane hurt herself.
the reflexive herself can only refer to the previously mentioned antecedent, in this case Jane, and can never refer to a different female person. In this example, the variable herself is bound to the noun Jane that occurs in subject position. Indicating the coindexation, the first interpretation with Jane and herself coindexed is permissible, but the other interpretation where they are not coindexed is ungrammatical:
Janei hurt herselfi. (interpretation #1: herself = Jane)
*Janei hurt herselfj. (interpretation #2: herself = a female that is not Jane)
The coreference binding can be represented using a lambda expression as mentioned in the previous Formal explanation section. The sentence with the reflexive could be represented as
(λx.x hurt x)Jane
in which Jane is the subject referent argument and λx.x hurt x is the predicate function (a lambda abstraction) with the lambda notation and x indicating both the semantic subject and the semantic object of sentence as being bound. This returns the semantic interpretation JANE hurt JANE with JANE being the same person.
Pronouns can also behave in a different way. In the sentence below
Ashley hit her.
the pronoun her can only refer to a female that is not Ashley. This means that it can never have a reflexive meaning equivalent to Ashley hit herself. The grammatical and ungrammatical interpretations are:
*Ashleyi hit heri. (interpretation #1: her = Ashley)
Ashleyi hit herj. (interpretation #2: her = a female that is not Ashley)
The first interpretation is impossible. Only the second interpretation is permitted by the grammar.
Thus, it can be seen that reflexives and reciprocals are bound variables (known technically as anaphors) while true pronouns are free variables in some grammatical structures but variables that cannot be bound in other grammatical structures. The binding phenomena found in natural languages was particularly important to the syntactic government and binding theory (see also: Binding (linguistics)).
| Mathematics | Mathematical logic | null |
147484 | https://en.wikipedia.org/wiki/Plant%20pathology | Plant pathology | Plant pathology or phytopathology is the scientific study of plant diseases caused by pathogens (infectious organisms) and environmental conditions (physiological factors). Plant pathology involves the study of pathogen identification, disease etiology, disease cycles, economic impact, plant disease epidemiology, plant disease resistance, how plant diseases affect humans and animals, pathosystem genetics, and management of plant diseases.
Plant pathogenicity
Plant pathogens, organisms that cause infectious plant diseases, include fungi, oomycetes, bacteria, viruses, viroids, virus-like organisms, phytoplasmas, protozoa, nematodes and parasitic plants.
In most plant pathosystems, virulence depends on hydrolases and enzymes that degrade the cell wall. The vast majority of these act on pectins (for example, pectinesterase, pectate lyase, and pectinases). For microbes, the cell wall polysaccharides are both a food source and a barrier to be overcome. Many pathogens grow opportunistically when the host breaks down its own cell walls, most often during fruit ripening. Unlike human and animal pathology, plant pathology usually focuses on a single causal organism; however, some plant diseases have been shown to be interactions between multiple pathogens.
To colonize a plant, pathogens have specific pathogenicity factors, of five main types: uses of cell wall–degrading enzymes, toxins, effector proteins, phytohormones, and exopolysaccharides.
Cell wall-degrading enzymes: These are used to break down the plant cell wall in order to release the nutrients inside and include esterases, glycosyl hydrolases, lyases and oxidoreductases.
Toxins: These can be non-host-specific, which damage all plants, or host-specific, which cause damage only on a host plant.
Effector proteins: These can be secreted by pathogens such as bacteria, fungi, and oomycetes into the extracellular environment or directly into the host cell, often via the Type three secretion system. Some effectors are known to suppress host defense processes. This can include reducing the plant's internal signaling mechanisms or reduction of phytochemicals production.
Phytohormones are chemicals used by plants for signaling; pathogens can produce these to modify plant growth to their own advantage.
Exopolysaccharides are mostly small chains of sugars that help pathogens to adhere to a plant's surface, enabling them to begin the process of infection.
Physiological plant disorders
Some abiotic disorders can be confused with pathogen-induced disorders. Abiotic causes include natural processes such as drought, frost, snow and hail; flooding and poor drainage; nutrient deficiency; deposition of mineral salts such as sodium chloride and gypsum; windburn and breakage by storms; and wildfires.
Epidemiology
Epidemiology is the study of factors affecting the outbreak and spread of infectious diseases.
A disease triangle describes the basic factors required for plant diseases. These are the host plant, the pathogen, and the environment. Any one of these can be modified to control a disease.
Disease resistance
Plant disease resistance is the ability of a plant to prevent and terminate infections from plant pathogens. Structures that help plants prevent pathogens from entering are the cuticular layer, cell walls, and stomata guard cells. Once pathogens have overcome these barriers, plant receptors initiate signaling pathways to create molecules to compete against the foreign molecules. These pathways are influenced and triggered by genes within the host plant and can manipulated by genetic breeding to create resistant varieties.
Management
Detection
Ancient methods of leaf examination and breaking open plant material by hand are now augmented by newer technologies. These include molecular pathology assays such as polymerase chain reaction (PCR), RT-PCR and loop-mediated isothermal amplification (LAMP). Although PCR can detect multiple molecular targets in a single solution there are limits. Bertolini et al. 2001, Ito et al. 2002, and Ragozzino et al. 2004 developed PCR methods for multiplexing six or seven plant pathogen molecular products and Persson et al. 2005 for multiplexing four with RT-PCR. More extensive molecular diagnosis requires PCR arrays. The primary detection method used worldwide is enzyme linked immunosorbent assay.
Biological
Crop rotation is a traditional and sometimes effective means of preventing a parasitic population from becoming well-established. For example, protection against infection by Agrobacterium tumefaciens, which causes gall diseases in many plants, by dipping cuttings in suspensions of Agrobacterium radiobacter before inserting them in the ground to take root.
History
Plant pathology has developed from antiquity, starting with Theophrastus in the ancient era, but scientific study began in the Early Modern period with the invention of the microscope, and developed in the 19th century.
Notable People in Plant Pathology
George Washington Carver
Anton de Bary
Erwin Frink Smith
Agnes Robertson Arber
| Biology and health sciences | Botany | Biology |
147518 | https://en.wikipedia.org/wiki/Ladder | Ladder | A ladder is a vertical or inclined set of rungs or steps commonly used for climbing or descending. There are two types: rigid ladders that are self-supporting or that may be leaned against a vertical surface such as a wall, and rollable ladders, such as those made of rope or aluminium, that may be hung from the top. The vertical members of a rigid ladder are called stringers or rails (US) or stiles (UK). Rigid ladders are usually portable, but some types are permanently fixed to a structure, building, or equipment. They are commonly made of metal, wood, or fiberglass, but they have been known to be made of tough plastic.
Historical usages
Ladders are ancient tools and technology. A ladder is featured in a Mesolithic rock painting that is at least 10,000 years old, depicted in the Spider Caves in Valencia, Spain. The painting depicts two humans using a ladder to reach a wild honeybee nest to harvest honey. The ladder is depicted as long and flexible, possibly made out of some sort of grass.
Variations
Rigid ladders
Rigid ladders are available in many forms, such as:
Accommodation ladder, portable steps down the side of a ship for boarding.
Assault ladder, used in siege warfare to assist in climbing walls and crossing moats.
Attic ladder, pulled down from the ceiling to allow access to an attic or loft.
, a ladder laid horizontally to act as a passage between two points separated by a drop.
Boarding ladder, a ladder used to climb onto a vehicle. May be rigid or flexible, also boarding step(s), and swim ladder.
Cat ladder (US chicken ladder), a lightweight ladder frame used on steep roofs to prevent workers from sliding.
Chicken ladder, a ladder comprising a single central stile with each rung projecting on either side and used by chickens to climb into a coop.
Christmas tree ladder, a type of boarding ladder for divers which has a single central rail and is open at the sides to allow the diver to climb the ladder while wearing swimfins.
, a fixed ladder with a lower sliding part. A system of counterweights is used to let the lower sliding part descend gently when released.
Extension ladder or "telescopic ladder", a fixed ladder divided into two or more lengths for more convenient storage; the lengths can be slid together for storage or slid apart to expand the length of the ladder; a pulley system may be fitted so that the ladder can be easily extended by an operator on the ground then locked in place using the dogs and pawls. 65 ft (20 m), 50 ft (15 m) and some 35 ft (10 m) extension ladders for fire service use "bangor poles", "tormentor poles" or "stay poles" to help raise, pivot, steady, extend, place, retract and lower them due to the heavy weight.
Fixed ladder, two side members joined by several rungs; affixed to structure with no moving parts.
, a ladder in the step ladder style with one or more (usually no more than three) one-way hinges. Ideal for use on uneven ground (e.g. stairs), as a trestle or when fully extended a Fixed ladder. Some variations feature a central one-way hinge with extensible locking legs.
Hook ladder or pompier ladder, a rigid ladder with a hook at the top to grip a windowsill; used by firefighters.
Mobile Safety Steps are self-supporting structures that have wheels or castors making them easy to move. They sometimes have a small upper platform and a hand rail to assist in moving up and down the steps.
Orchard ladder, a three legged step ladder with the third leg made so that it can be inserted between tree branches for fruit picking.
, a step ladder with a large platform area and a top handrail for the user to hold while working on the platform.
, a ladder that looks like a drainpipe but can be deployed instantly when required.
Roof ladder, a rigid ladder with a large hook at the top to grip the ridge of a pitched roof.
, also known as a builder's ladder, has sections that come apart and are interchangeable so that any number of sections can be connected.
Step ladder, a self-supporting portable ladder hinged in the middle to form an inverted V, with stays to keep the two halves at a fixed angle. Step ladders have flat steps and a hinged back.
Swim ladder, a ladder used by swimmers to get out of the water, often on boats.
, commonly used to refer to a hybrid between a step ladder and an extension ladder with 360-degree hinges; has three parts and can be taken apart to form two step ladders; e.g. Little Giant.
, an "A-Frame"-style ladder with a telescoping center section.
Turntable ladder, an extension ladder fitted to rotating platform on top of a fire truck.
Vertically rising ladder, designed to climb high points and facilitate suspending at said high points.
X-deck ladder, a US patented ladder design that is a combination ladder and scaffold.
Rigid ladders were originally made of wood, but in the 20th century aluminium became more common because of its lighter weight. Ladders with fiberglass stiles are used for working on or near overhead electrical wires, because fiberglass is an electrical insulator. Henry Quackenbush patented the extension ladder in 1867.
Flexible ladders
Rope ladders or Jacob's ladders are used where storage space is extremely limited, weight must be kept to a minimum, or in instances where the object to be climbed is too curved to use a rigid ladder. They may have rigid or flexible rungs. Climbing a rope ladder requires more skill than climbing a rigid ladder, because the ladder tends to swing like a pendulum. Jacob's ladders used on a ship are used mostly for emergencies or for temporary access to the side of a ship. Steel and aluminum wire ladders are sometimes used in vertical caving, having developed from rope ladders with wooden rungs. Flexible ladders are also sometimes used as swim ladders on boats.
Uses
Dissipative ladders are portable ladders built to ESD (Electrostatic Discharge) standard. Electrostatic Discharge is a natural occurrence in which electricity is passed through the body, or other conductors, and discharges onto some object. For example, the shock sometimes felt when a doorknob is touched is an ESD. This natural occurrence is a very important topic in the field of electronics assembly due to the costly damage ESDs can cause to sensitive electronic equipment. Dissipative ladders are ladders with controlled electrical resistance: the resistance slows the transfer of charge from one point to another, offering increased protection during ESD events: ≥105 and < 1012 Ω / square.
Boarding and pool ladders, also swim ladders and dive ladders. A ladder may be used on the side or stern of a boat, to climb into it from the water, and in a swimming pool, to climb out and sometimes in. Swimming pool ladders are usually made from plastic, wood or metal steps with a textured upper surface for grip and metal rails at the sides to support the steps and as handrails for the user, and are usually fixed in place. Boarding ladders for boats may be fixed, but are usually portable, and often fold away when not in use to avoid drag when under way. Boarding ladders may also be used for other types of vehicle, or boarding steps which are supported directly by the vehicle structure.
Assault ladders
Safety
The most common injury made by ladder climbers is bruising from falling off a ladder, but bone fractures are common and head injuries are also likely, depending on the nature of the accident. Ladders can cause injury if they slip on the ground and fall. To avoid this, they tend to have plastic feet or base pads which increase friction with the ground. However, if the plastic is badly worn, the aluminium may contact the ground increasing the chance of an accident. Ladder stabilizers are also available to increase the ladder's grip on the ground. One of the first ladder stabilizers or ladder feet was offered in 1936 and today they are standard equipment on most large ladders.
A ladder standoff, or stay, is a device fitted to the top of a ladder to hold it away from the wall. This enables the ladder to clear overhanging obstacles, such as the eaves of a roof, and increases the safe working height for a given length of ladder because of the increased separation distance of the two contact points at the top of the ladder.
It has become increasingly common to provide anchor points on buildings to which the top rung of an extension ladder can be attached, especially for activities like window cleaning, especially if a fellow worker is not available for "footing" the ladder. Footing occurs when another worker stands on the lowest rung and so provides much greater stability to the ladder when being used. However footing a ladder should be seen as a last resort for a safe placement. The anchor point is usually a ring cemented into a slot in the brick wall to which the rungs of a ladder can be attached using rope for example, or a carabiner.
If a leaning ladder is placed at the wrong angle, the risk of a fall is greatly increased. The safest angle for a ladder is 75.5°; if it is too shallow, the bottom of the ladder is at risk of sliding, and if it is too steep, the ladder may fall backwards. This angle is achieved by following the 4 to 1 rule for a ladder placed on a vertical wall: for every four feet of vertical height, the ladder foot should move one foot from the wall. Both scenarios can cause significant injury, and are especially important in industries like construction, which require heavy use of ladders.
Ladder classes
The European Union and the United Kingdom established a ladder certification system – ladder classes – for any ladders manufactured or sold in Europe. The certification classes apply solely to ladders that are portable such as stepladders and extension ladders and are broken down into three types of certification. Each ladder certification is colour-coded to indicate the amount of weight the ladder is designed to hold, the certification class and its use. The color of the safety label specifies the class and use.
Class 1 ladder – for heavy-duty industrial uses, maximum load of 175 kg. Colour-coded blue to identify.
Class EN131 ladders – for commercial uses, maximum load of 150 kg. No specific colour code..
Class III ladders – for light, domestic uses, maximum load of 125 kg. Colour-coded red to identify.
Society and culture
A common superstition in English-speaking countries is that walking under a ladder is seen as bad luck. Some sources claim that this stems from the image of a ladder being propped up against a wall looking similar to a gallows, while others attribute it to ancient Egyptian traditions involving pyramids and triangles representing the trinity of the gods, and passing through the triangular shape made by a ladder against a wall was seen as desecration. Ladders have also been linked to the crucifixion of Christ, with author and scientist Charles Panati noting that many believe a ladder rested against the cross that Christ hung from, making it a symbol of wickedness, betrayal and death. In comedic children's media, the image of a character walking under a ladder being the cause or result of bad luck has become a common trope.
Image gallery
| Technology | Architectural elements | null |
147536 | https://en.wikipedia.org/wiki/Calcium%20oxide | Calcium oxide | Calcium oxide (formula: CaO), commonly known as quicklime or burnt lime, is a widely used chemical compound. It is a white, caustic, alkaline, crystalline solid at room temperature. The broadly used term lime connotes calcium-containing inorganic compounds, in which carbonates, oxides, and hydroxides of calcium, silicon, magnesium, aluminium, and iron predominate. By contrast, quicklime specifically applies to the single compound calcium oxide. Calcium oxide that survives processing without reacting in building products, such as cement, is called free lime.
Quicklime is relatively inexpensive. Both it and the chemical derivative calcium hydroxide (of which quicklime is the base anhydride) are important commodity chemicals.
Preparation
Calcium oxide is usually made by the thermal decomposition of materials, such as limestone or seashells, that contain calcium carbonate (CaCO3; mineral calcite) in a lime kiln. This is accomplished by heating the material to above , a process called calcination or lime-burning, to liberate a molecule of carbon dioxide (CO2), leaving quicklime behind. This is also one of the few chemical reactions known in prehistoric times.
CaCO3(s) → CaO(s) + CO2(g)
The quicklime is not stable and, when cooled, will spontaneously react with CO2 from the air until, after enough time, it will be completely converted back to calcium carbonate unless slaked with water to set as lime plaster or lime mortar.
Annual worldwide production of quicklime is around 283 million tonnes. China is by far the world's largest producer, with a total of around 170 million tonnes per year. The United States is the next largest, with around 20 million tonnes per year.
Approximately 1.8t of limestone is required per 1.0t of quicklime. Quicklime has a high affinity for water and is a more efficient desiccant than silica gel. The reaction of quicklime with water is associated with an increase in volume by a factor of at least 2.5.
Hydroxyapatite's free CaO content rises with increased calcination temperatures and longer times. It also pinpoints particular temperature cutoffs and durations that impact the production of CaO, offering information on how calcination parameters impact the composition of the material.
Uses
The major use of quicklime is in the basic oxygen steelmaking (BOS) process. Its usage varies from about per ton of steel. The quicklime neutralizes the acidic oxides, SiO2, Al2O3, and Fe2O3, to produce a basic molten slag.
Ground quicklime is used in the production of aerated concrete such as blocks with densities of ca. .
Quicklime and hydrated lime can considerably increase the load carrying capacity of clay-containing soils. They do this by reacting with finely divided silica and alumina to produce calcium silicates and aluminates, which possess cementing properties.
Small quantities of quicklime are used in other processes; e.g., the production of glass, calcium aluminate cement, and organic chemicals.
Heat: Quicklime releases thermal energy by the formation of the hydrate, calcium hydroxide, by the following equation:
CaO (s) + H2O (l) Ca(OH)2 (aq) (ΔHr = −63.7kJ/mol of CaO)
As it hydrates, an exothermic reaction results and the solid puffs up. The hydrate can be reconverted to quicklime by removing the water by heating it to redness to reverse the hydration reaction. One litre of water combines with approximately of quicklime to give calcium hydroxide plus 3.54 MJ of energy. This process can be used to provide a convenient portable source of heat, as for on-the-spot food warming in a self-heating can, cooking, and heating water without open flames. Several companies sell cooking kits using this heating method.
It is known as a food additive to the FAO as an acidity regulator, a flour treatment agent and as a leavener. It has E number E529.
Light: When quicklime is heated to , it emits an intense glow. This form of illumination is known as a limelight, and was used broadly in theatrical productions before the invention of electric lighting.
Cement: Calcium oxide is a key ingredient for the process of making cement.
As a cheap and widely available alkali. About 50% of the total quicklime production is converted to calcium hydroxide before use. Both quick- and hydrated lime are used in the treatment of drinking water.
Petroleum industry: Water detection pastes contain a mix of calcium oxide and phenolphthalein. Should this paste come into contact with water in a fuel storage tank, the CaO reacts with the water to form calcium hydroxide. Calcium hydroxide has a high enough pH to turn the phenolphthalein a vivid purplish-pink color, thus indicating the presence of water.
Chemical pulping: Calcium oxide is used to make calcium hydroxide, which is used to regenerate sodium hydroxide from sodium carbonate in the chemical recovery at kraft pulp mills.
Plaster: There is archeological evidence that Pre-Pottery Neolithic B humans used limestone-based plaster for flooring and other uses. Such Lime-ash floor remained in use until the late nineteenth century.
Chemical or power production: Solid sprays or slurries of calcium oxide can be used to remove sulfur dioxide from exhaust streams in a process called flue-gas desulfurization.
Carbon capture and storage: Calcium oxide can be used to capture carbon dioxide from flue gases in a process called calcium looping.
Mining: Compressed lime cartridges exploit the exothermic properties of quicklime to break rock. A shot hole is drilled into the rock in the usual way and a sealed cartridge of quicklime is placed within and tamped. A quantity of water is then injected into the cartridge and the resulting release of steam, together with the greater volume of the residual hydrated solid, breaks the rock apart. The method does not work if the rock is particularly hard.
Disposal of corpses: Historically, it was mistakenly believed that quicklime was efficacious in accelerating the decomposition of corpses. The application of quicklime can, in fact, promote preservation. Quicklime can aid in eradicating the stench of decomposition, which may have led people to the erroneous conclusion.
It has been determined that the durability of ancient Roman concrete is attributed in part to the use of quicklime as an ingredient. Combined with hot mixing, the quicklime creates macro-sized lime clasts with a characteristically brittle nano-particle architecture. As cracks form in the concrete, they preferentially pass through the structurally weaker lime clasts, fracturing them. When water enters these cracks it creates a calcium-saturated solution which can recrystallize as calcium carbonate, quickly filling the crack.
The thermochemical heat storage mechanism is greatly impacted by the sintering of CaO and CaCO3. It demonstrates that the storage materials become less reactive and denser at increasing temperatures. It also pinpoints particular sintering processes and variables influencing the efficiency of these materials in heat storage.
Weapon
In 80 BC, the Roman general Sertorius deployed choking clouds of caustic lime powder to defeat the Characitani of Hispania, who had taken refuge in inaccessible caves. A similar dust was used in China to quell an armed peasant revolt in 178 AD, when lime chariots equipped with bellows blew limestone powder into the crowds.
Quicklime is also thought to have been a component of Greek fire. Upon contact with water, quicklime would increase its temperature above and ignite the fuel.
David Hume, in his History of England, recounts that early in the reign of Henry III, the English Navy destroyed an invading French fleet by blinding the enemy fleet with quicklime. Quicklime may have been used in medieval naval warfare – up to the use of "lime-mortars" to throw it at the enemy ships.
Substitutes
Limestone is a substitute for lime in many applications, which include agriculture, fluxing, and sulfur removal. Limestone, which contains less reactive material, is slower to react and may have other disadvantages compared with lime, depending on the application; however, limestone is considerably less expensive than lime. Calcined gypsum is an alternative material in industrial plasters and mortars. Cement, cement kiln dust, fly ash, and lime kiln dust are potential substitutes for some construction uses of lime. Magnesium hydroxide is a substitute for lime in pH control, and magnesium oxide is a substitute for dolomitic lime as a flux in steelmaking.
Safety
Because of vigorous reaction of quicklime with water, quicklime causes severe irritation when inhaled or placed in contact with moist skin or eyes. Inhalation may cause coughing, sneezing, and labored breathing. It may then evolve into burns with perforation of the nasal septum, abdominal pain, nausea and vomiting. Although quicklime is not considered a fire hazard, its reaction with water can release enough heat to ignite combustible materials.
Mineral
Calcium oxide is also a separate mineral species (with the unit formula CaO), named 'Lime'. It has an isometric crystal system, and can form a solid solution series with monteponite. The crystal is brittle, pyrometamorphic, and is unstable in moist air, quickly turning into portlandite (Ca(OH)2).
| Physical sciences | Alkali oxide salts | Chemistry |
147566 | https://en.wikipedia.org/wiki/Otto%20cycle | Otto cycle | An Otto cycle is an idealized thermodynamic cycle that describes the functioning of a typical spark ignition piston engine. It is the thermodynamic cycle most commonly found in automobile engines.
The Otto cycle is a description of what happens to a gas as it is subjected to changes of pressure, temperature, volume, addition of heat, and removal of heat. The gas that is subjected to those changes is called the system. The system, in this case, is defined to be the fluid (gas) within the cylinder. Conversely, by describing the changes that take place within the system it also describes the system's effect on the environment. The purpose of the Otto cycle is to study the production of net work from the system that can propel a vehicle and its occupants in the environment.
The Otto cycle is constructed from:
Top and bottom of the loop: a pair of quasi-parallel and isentropic processes (frictionless, adiabatic reversible).
Left and right sides of the loop: a pair of parallel isochoric processes (constant volume).
The isentropic process of compression or expansion implies that there will be no inefficiency (loss of mechanical energy), and there be no transfer of heat into or out of the system during that process. The cylinder and piston are assumed to be impermeable to heat during that time. Work is performed on the system during the lower isentropic compression process. Heat flows into the Otto cycle through the left pressurizing process and some of it flows back out through the right depressurizing process. The summation of the work added to the system plus the heat added minus the heat removed yields the net mechanical work generated by the system.
Processes
The processes are described by:
Process 0–1 a mass of air is drawn into piston/cylinder arrangement at constant pressure.
Process 1–2 is an adiabatic (isentropic) compression of the charge as the piston moves from bottom dead center (BDC) to top dead center (TDC).
Process 2–3 is a constant-volume heat transfer to the working gas from an external source while the piston is at top dead center. This process is intended to represent the ignition of the fuel-air mixture and the subsequent rapid burning.
Process 3–4 is an adiabatic (isentropic) expansion (power stroke).
Process 4–1 completes the cycle by a constant-volume process in which heat is rejected from the air while the piston is at bottom dead center.
Process 1–0 the mass of air is released to the atmosphere in a constant pressure process.
The Otto cycle consists of isentropic compression, heat addition at constant volume, isentropic expansion, and rejection of heat at constant volume. In the case of a four-stroke Otto cycle, technically there are two additional processes: one for the exhaust of waste heat and combustion products at constant pressure (isobaric), and one for the intake of cool oxygen-rich air also at constant pressure; however, these are often omitted in a simplified analysis. Even though those two processes are critical to the functioning of a real engine, wherein the details of heat transfer and combustion chemistry are relevant, for the simplified analysis of the thermodynamic cycle, it is more convenient to assume that all of the waste-heat is removed during a single volume change.
History
The four-stroke engine was first patented by Alphonse Beau de Rochas in 1861. Before, in about 1854–57, two Italians (Eugenio Barsanti and Felice Matteucci) invented an engine that was rumored to be very similar, but the patent was lost.
The first person to build a working four-stroke engine, a stationary engine using a coal gas-air mixture for fuel (a gas engine), was German engineer Nicolaus Otto. This is why the four-stroke principle today is commonly known as the Otto cycle and four-stroke engines using spark plugs often are called Otto engines.
Processes
The cycle has four parts: a mass containing a mixture of fuel and oxygen is drawn into the cylinder by the descending piston, it is compressed by the piston rising, the mass is ignited by a spark releasing energy in the form of heat, the resulting gas is allowed to expand as it pushes the piston down, and finally the mass is exhausted as the piston rises a second time. As the piston is capable of moving along the cylinder, the volume of the gas changes with its position in the cylinder. The compression and expansion processes induced on the gas by the movement of the piston are idealized as reversible, i.e., no useful work is lost through turbulence or friction and no heat is transferred to or from the gas during those two processes. After the expansion is completed in the cylinder, the remaining heat is extracted and finally the gas is exhausted to the environment. Mechanical work is produced during the expansion process and some of that used to compress the air mass of the next cycle. The mechanical work produced minus that used for the compression process is the net work gained and that can be used for propulsion or for driving other machines. Alternatively the net work gained is the difference between the heat produced and the heat removed.
Process 0–1 intake stroke (blue shade)
A mass of air (working fluid) is drawn into the cylinder, from 0 to 1, at atmospheric pressure (constant pressure) through the open intake valve, while the exhaust valve is closed during this process. The intake valve closes at point 1.
Process 1–2 compression stroke (B on diagrams)
Piston moves from crank end (BDC, bottom dead centre and maximum volume) to cylinder head end (TDC, top dead centre and minimum volume) as the working gas with initial state 1 is compressed isentropically to state point 2, through compression ratio . Mechanically this is the isentropic compression of the air/fuel mixture in the cylinder, also known as the compression stroke. This isentropic process assumes that no mechanical energy is lost due to friction and no heat is transferred to or from the gas, hence the process is reversible. The compression process requires that mechanical work be added to the working gas. Generally the compression ratio is around 9–10:1 for a typical engine.
Process 2–3 ignition phase (C on diagrams)
The piston is momentarily at rest at TDC. During this instant, which is known as the ignition phase, the air/fuel mixture remains in a small volume at the top of the compression stroke. Heat is added to the working fluid by the combustion of the injected fuel, with the volume essentially being held constant. The pressure rises and the ratio is called the "explosion ratio".
Process 3–4 expansion stroke (D on diagrams)
The increased high pressure exerts a force on the piston and pushes it towards the BDC. Expansion of working fluid takes place isentropically and work is done by the system on the piston. The volume ratio is called the "isentropic expansion ratio". (For the Otto cycle is the same as the compression ratio ). Mechanically this is the expansion of the hot gaseous mixture in the cylinder known as expansion (power) stroke.
Process 4–1 idealized heat rejection (A on diagrams)
The piston is momentarily at rest at BDC. The working gas pressure drops instantaneously from point 4 to point 1 during a constant volume process as heat is removed to an idealized external sink that is brought into contact with the cylinder head. In modern internal combustion engines, the heat-sink may be surrounding air (for low powered engines), or a circulating fluid, such as coolant. The gas has returned to state 1.
Process 1–0 exhaust stroke
The exhaust valve opens at point 1. As the piston moves from "BDC" (point 1) to "TDC" (point 0) with the exhaust valve opened, the gaseous mixture is vented to the atmosphere and the process starts anew.
Cycle analysis
In this process 1–2 the piston does work on the gas and in process 3–4 the gas does work on the piston during those isentropic compression and expansion processes, respectively. Processes 2–3 and 4–1 are isochoric processes; heat is transferred into the system from 2—3 and out of the system from 4–1 but no work is done on the system or extracted from the system during those processes. No work is done during an isochoric (constant volume) process because addition or removal of work from a system requires the movement of the boundaries of the system; hence, as the cylinder volume does not change, no shaft work is added to or removed from the system.
Four different equations are used to describe those four processes. A simplification is made by assuming changes of the kinetic and potential energy that take place in the system (mass of gas) can be neglected and then applying the first law of thermodynamics (energy conservation) to the mass of gas as it changes state as characterized by the gas's temperature, pressure, and volume.
During a complete cycle, the gas returns to its original state of temperature, pressure and volume, hence the net internal energy change of the system (gas) is zero. As a result, the energy (heat or work) added to the system must be offset by energy (heat or work) that leaves the system. In the analysis of thermodynamic systems, the convention is to account energy that enters the system as positive and energy that leaves the system is accounted as negative.
Equation 1a.
During a complete cycle, the net change of energy of the system is zero:
The above states that the system (the mass of gas) returns to the original thermodynamic state it was in at the start of the cycle.
Where is energy added to the system from 1–2–3 and is energy removed from the system from 3–4–1. In terms of work and heat added to the system
Equation 1b:
Each term of the equation can be expressed in terms of the internal energy of the gas at each point in the process:
The energy balance Equation 1b becomes
To illustrate the example we choose some values to the points in the illustration:
These values are arbitrarily but rationally selected. The work and heat terms can then be calculated.
The energy added to the system as work during the compression from 1 to 2 is
The energy added to the system as heat from point 2 to 3 is
The energy removed from the system as work during the expansion from 3 to 4 is
The energy removed from the system as heat from point 4 to 1 is
The energy balance is
Note that energy added to the system is counted as positive and energy leaving the system is counted as negative and the summation is zero as expected for a complete cycle that returns the system to its original state.
From the energy balance the work out of the system is:
The net energy out of the system as work is -1, meaning the system has produced one net unit of energy that leaves the system in the form of work.
The net heat out of the system is:
As energy added to the system as heat is positive. From the above it appears as if the system gained one unit of heat. This matches the energy produced by the system as work out of the system.
Thermal efficiency is the quotient of the net work from the system, to the heat added to system.
Equation 2:
Alternatively, thermal efficiency can be derived by strictly heat added and heat rejected.
Supplying the fictitious values
In the Otto cycle, there is no heat transfer during the process 1–2 and 3–4 as they are isentropic processes. Heat is supplied only during the constant volume processes 2–3 and heat is rejected only during the constant volume processes 4–1.
The above values are absolute values that might, for instance , have units of joules (assuming the MKS system of units are to be used) and would be of use for a particular engine with particular dimensions. In the study of thermodynamic systems the extensive quantities such as energy, volume, or entropy (versus intensive quantities of temperature and pressure) are placed on a unit mass basis, and so too are the calculations, making those more general and therefore of more general use. Hence, each term involving an extensive quantity could be divided by the mass, giving the terms units of joules/kg (specific energy), meters3/kg (specific volume), or joules/(kelvin·kg) (specific entropy, heat capacity) etc. and would be represented using lower case letters, u, v, s, etc.
Equation 1 can now be related to the specific heat equation for constant volume. The specific heats are particularly useful for thermodynamic calculations involving the ideal gas model.
Rearranging yields:
Inserting the specific heat equation into the thermal efficiency equation (Equation 2) yields.
Upon rearrangement:
Next, noting from the diagrams (see isentropic relations for an ideal gas), thus both of these can be omitted. The equation then reduces to:
Equation 2:
Since the Otto cycle uses isentropic processes during the compression (process 1 to 2) and expansion (process 3 to 4) the isentropic equations of ideal gases and the constant pressure/volume relations can be used to yield Equations 3 & 4.
Equation 3:
Equation 4:
where
is the specific heat ratio
The derivation of the previous equations are found by solving these four equations respectively (where is the specific gas constant):
Further simplifying Equation 4, where is the compression ratio :
Equation 5:
From inverting Equation 4 and inserting it into Equation 2 the final thermal efficiency can be expressed as:
Equation 6:
From analyzing equation 6 it is evident that the Otto cycle efficiency depends directly upon the compression ratio . Since the for air is 1.4, an increase in will produce an increase in . However, the for combustion products of the fuel/air mixture is often taken at approximately 1.3.
The foregoing discussion implies that it is more efficient to have a high compression ratio. The standard ratio is approximately 10:1 for typical automobiles. Usually this does not increase much because of the possibility of autoignition, or "knock", which places an upper limit on the compression ratio. During the compression process 1–2 the temperature rises, therefore an increase in the compression ratio causes an increase in temperature. Autoignition occurs when the temperature of the fuel/air mixture becomes too high before it is ignited by the flame front. The compression stroke is intended to compress the products before the flame ignites the mixture. If the compression ratio is increased, the mixture may auto-ignite before the compression stroke is complete, leading to "engine knocking". This can damage engine components and will decrease the brake horsepower of the engine.
Power
The power produced by the Otto cycle is an energy developed per unit of time. The Otto engines are called four-stroke engines.
The intake stroke and compression stroke require one rotation of the engine crankshaft. The power stroke and exhaust stroke require another rotation. For two rotations there is one work generating stroke..
From the above cycle analysis the net work produced by the system :
(again, using the sign convention, the minus sign implies energy is leaving the system as work)
If the units used were MKS the cycle would have produced one joule of energy in the form of work. For an engine of a particular displacement, such as one liter, the mass of gas of the system can be calculated assuming the engine is operating at standard temperature (20 °C) and pressure (1 atm). Using the Universal Gas Law the mass of one liter of gas is at room temperature and sea level pressure:
V=0.001 m3, R=0.286 kJ/(kg·K), T=293 K, P=101.3 kN/m2
M=0.00121 kg
At an engine speed of 3000 RPM there are 1500 work-strokes/minute or 25 work-strokes/second.
Power is 25 times that since there are 25 work-strokes/second
If the engine uses multiple cylinders with the same displacement, the result would be multiplied by the number of cylinders. These results are the product of the values of the internal energy that were assumed for the four states of the system at the end each of the four strokes (two rotations). They were selected only for the sake of illustration, and are obviously of low value. Substitution of actual values from an actual engine would produce results closer to that of the engine. Whose results would be higher than the actual engine as there are many simplifying assumptions made in the analysis that overlook inefficiencies. Such results would overestimate the power output.
Increasing power and efficiency
The difference between the exhaust and intake pressures and temperatures means that some increase in efficiency can be gained by use of a turbocharger, removing from the exhaust flow some part of the remaining energy and transferring that to the intake flow to increase the intake pressure. A gas turbine can extract useful work energy from the exhaust stream and that can then be used to pressurize the intake air. The pressure and temperature of the exhausting gases would be reduced as they expand through the gas turbine and that work is then applied to the intake gas stream, increasing its pressure and temperature. The transfer of energy amounts to an efficiency improvement and the resulting power density of the engine is also improved. The intake air is typically cooled so as to reduce its volume as the work produced per stroke is a direct function of the amount of mass taken into the cylinder; denser air will produce more work per cycle. Practically speaking the intake air mass temperature must also be reduced to prevent premature ignition in a petrol fueled engine; hence, an intercooler is used to remove some energy as heat and so reduce the intake temperature. Such a scheme both increases the engine's efficiency and power.
The application of a supercharger driven by the crankshaft does increase the power output (power density) but does not increase efficiency as it uses some of the net work produced by the engine to pressurize the intake air and fails to extract otherwise wasted energy associated with the flow of exhaust at high temperature and a pressure to the ambient.
| Physical sciences | Thermodynamics | Physics |
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