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147572 | https://en.wikipedia.org/wiki/Camshaft | Camshaft | A camshaft is a shaft that contains a row of pointed cams in order to convert rotational motion to reciprocating motion. Camshafts are used in piston engines (to operate the intake and exhaust valves), mechanically controlled ignition systems and early electric motor speed controllers.
Camshafts in piston engines are usually made from steel or cast iron, and the shape of the cams greatly affects the engine's characteristics.
History
Trip hammers are one of the early uses of a form of cam to convert rotating motion, e.g. from a waterwheel, into the reciprocating motion of a hammer used in forging or to pound grain. Evidence for these exists back to the Han dynasty in China, and they were widespread by the medieval period.
Once the rotative version of the steam engine was developed in the late 18th century, the operation of the valve gear was usually by an eccentric, which turned the rotation of the crankshaft into reciprocating motion of the valve gear, normally a slide valve. Camshafts more like those seen later in internal combustion engines were used in some steam engines, most commonly where high pressure steam (such as that generated from a flash steam boiler), required the use of poppet valves, or piston valves. For examples see the Uniflow steam engine, and the Gardner-Serpollet steam cars, which also included axially sliding the camshaft to achieve variable valve timing.
Among the first cars to utilize engines with single overhead camshafts were the Maudslay, designed by Alexander Craig and introduced in 1902 and the Marr Auto Car designed by Michigan native Walter Lorenzo Marr in 1903.
Piston engines
In piston engines, the camshaft is used to operate the intake and exhaust valves. The camshaft consists of a cylindrical rod running the length of the cylinder bank with a number of cams (discs with protruding cam lobes) along its length, one for each valve. As the cam rotates, the lobe presses on the valve (or an intermediate mechanism), thus pushing it open. Typically, a valve spring is used to push the valve in the opposite direction, thus closing the valve once the cam rotates past the highest point of its lobe.
Construction
Camshafts are made from metal and are usually solid, although hollow camshafts are sometimes used. The materials used for a camshaft are usually either:
Cast iron: Commonly used in high volume production, chilled iron camshafts have good wear resistance since the chilling process hardens them.
Billet steel: For high-performance engines or camshafts produced in small quantities, steel billet is sometimes used. This is a much more time-consuming process, and is generally more expensive than other methods. The method of construction is usually either forging, machining, casting or hydroforming.
Location in engine
Many early internal combustion engines used a cam-in-block layout (such flathead, IOE or T-head layouts), whereby the camshaft is located within the engine block near the bottom of the engine. Early flathead engines locate the valves in the block and the cam acts directly on those valves. In an overhead valve engine, which came later, the cam follower presses on a pushrod which transfers the motion to the top of the engine, where a rocker opens the intake/exhaust valve. Although largely replaced by SOHC and DOHC layouts in modern automobile engines, the older overhead valve layout is still used in many industrial engines, due to its smaller size and lower cost.
As engine speeds increased through the 20th century, single overhead camshaft (SOHC) engines— where the camshaft is located within the cylinder head near the top of the engine— became increasingly common, followed by double overhead camshaft (DOHC) engines in more recent years. For OHC and DOHC engines, the camshaft operates the valve directly or via a short rocker arm.
The valvetrain layout is defined according to the number of camshafts per cylinder bank. Therefore, a V6 engine with a total of four camshafts - two camshafts per cylinder bank - is usually referred to as a double overhead camshaft engine (although colloquially they are sometimes referred to as "quad-cam" engines).
Drive systems
Accurate control of the position and speed of the camshaft is critically important in allowing the engine to operate correctly. The camshaft is usually driven either directly, via a toothed rubber "timing belt"' or via a steel roller "timing chain". Gears have also occasionally been used to drive the camshaft. In some designs the camshaft also drives the distributor, oil pump, fuel pump and occasionally the power steering pump.
Alternative drive systems used in the past include a vertical shaft with bevel gears at each end (e.g. pre-World War I Peugeot and Mercedes Grand Prix Cars and the Kawasaki W800 motorcycle) or a triple eccentric with connecting rods (e.g. the Leyland Eight car).
In a two-stroke engine that uses a camshaft, each valve is opened once for every rotation of the crankshaft; in these engines, the camshaft rotates at the same speed as the crankshaft. In a four-stroke engine, the valves are opened only half as often, therefore the camshaft is geared to rotate at half the speed of the crankshaft.
Performance characteristics
Duration
The camshaft's duration determines how long the intake/exhaust valve is open for, therefore it is a key factor in the amount of power that an engine produces. A longer duration can increase power at high engine speeds (RPM), however this can come with the trade-off of less torque being produced at low RPM.
The duration measurement for a camshaft is affected by the amount of lift that is chosen as the start and finish point of the measurement. A lift value of is often used as a standard measurement procedure, since this is considered most representative of the lift range that defines the RPM range in which the engine produces peak power. The power and idle characteristics of a camshaft with the same duration rating that has been determined using different lift points (for example 0.006 or 0.002 inches) could be much different to a camshaft with a duration rated using lift points of 0.05 inches.
A secondary effect of increased duration can be increased overlap, which determines the length of time that both the intake and exhaust valves are open. It is overlap which most affects idle quality, in as much as the "blow-through" of the intake charge immediately back out through the exhaust valve which occurs during overlap reduces engine efficiency, and is greatest during low RPM operation. In general, increasing a camshaft's duration typically increases the overlap, unless the Lobe Separation Angle is increased to compensate.
A lay person can readily spot a long duration camshaft by observing the broad surface of the lobe where the cam pushes the valve open for a large number of degrees of crankshaft rotation. This will be visibly greater than the more pointed camshaft lobe bump that is observed on lower duration camshafts.
Lift
The camshaft's lift determines the distance between the valve and the valve seat (i.e. how far open the valve is). The farther the valve rises from its seat the more airflow can be provided, thus increasing the power produced. Higher valve lift can have the same effect of increasing peak power as increased duration, without the downsides caused by increased valve overlap. Most overhead valve engines have a rocker ratio of greater than one, therefore the distance that the valve opens (the valve lift) is greater than the distance from the peak of the camshaft's lobe to the base circle (the camshaft lift).
There are several factors which limit the maximum amount of lift possible for a given engine. Firstly, increasing lift brings the valves closer to the piston, so excessive lift could cause the valves to get struck and damaged by the piston. Secondly, increased lift means a steeper camshaft profile is required, which increases the forces needed to open the valve. A related issue is valve float at high RPM, where the spring tension does not provide sufficient force to either keep the valve following the cam at its apex or prevent the valve from bouncing when it returns to the valve seat. This could be a result of a very steep rise of the lobe, where the cam follower separates from the cam lobe (due to the valvetrain inertia being greater than the closing force of the valve spring), leaving the valve open for longer than intended. Valve float causes a loss of power at high RPM and in extreme situations can result in a bent valve if it gets struck by the piston.
Timing
The timing (phase angle) of the camshaft relative to the crankshaft can be adjusted to shift an engine's power band to a different RPM range. Advancing the camshaft (shifting it to ahead of the crankshaft timing) increases low RPM torque, while retarding the camshaft (shifting it to after the crankshaft) increases high RPM power. The required changes are relatively small, often in the order of 5 degrees.
Modern engines which have variable valve timing are often able to adjust the timing of the camshaft to suit the RPM of the engine at any given time. This avoids the above compromise required when choosing a fixed cam timing for use at both high and low RPM.
Lobe separation angle
The lobe separation angle (LSA, also called lobe centreline angle) is the angle between the centreline of the intake lobes and the centreline of the exhaust lobes. A higher LSA reduces overlap, which improves idle quality and intake vacuum, however using a wider LSA to compensate for excessive duration can reduce power and torque outputs. In general, the optimal LSA for a given engine is related to the ratio of the cylinder volume to intake valve area.
Functionality
Camshafts are integral components of internal combustion engines, responsible for controlling the opening and closing of the engine's intake and exhaust valves. As the camshaft rotates, its lobes push against the valves, allowing the intake of air and fuel and the expulsion of exhaust gases. This synchronized process is crucial for optimizing engine performance, fuel efficiency, and emissions control. Without precisely engineered camshafts, the smooth and efficient operation of an engine would be compromised.
Alternatives
The most common methods of valve actuation involve camshafts and valve springs, however alternate systems have occasionally been used on internal combustion engines:
Desmodromic valves, where the valves are positively closed by a cam and leverage system rather than springs. This system has been used on various Ducati racing and road motorcycles since it was introduced on the 1956 Ducati 125 Desmo racing bike.
Camless piston engine, which use electromagnetic, hydraulic, or pneumatic actuators. First used in turbocharged Renault Formula 1 engines in the mid-1980s and slated for road car use in the Koenigsegg Gemera.
Wankel engine, a rotary engine which uses neither pistons nor valves. Most notably used by Mazda from the 1967 Mazda Cosmo until the Mazda RX-8 was discontinued in 2012.
Electric motor speed controllers
Before the advent of solid state electronics, camshaft controllers were used to control the speed of electric motors. A camshaft, driven by an electric motor or a pneumatic motor, was used to operate contactors in sequence. By this means, resistors or tap changers were switched in or out of the circuit to vary the speed of the main motor. This system was mainly used in electric train motors (i.e. EMUs and locomotives).
| Technology | Mechanisms | null |
147685 | https://en.wikipedia.org/wiki/Geomagnetic%20storm | Geomagnetic storm | A geomagnetic storm, also known as a magnetic storm, is a temporary disturbance of the Earth's magnetosphere caused by a solar wind shock wave.
The disturbance that drives the magnetic storm may be a solar coronal mass ejection (CME) or (much less severely) a corotating interaction region (CIR), a high-speed stream of solar wind originating from a coronal hole. The frequency of geomagnetic storms increases and decreases with the sunspot cycle. During solar maxima, geomagnetic storms occur more often, with the majority driven by CMEs.
The increase in the solar wind pressure initially compresses the magnetosphere. The solar wind's magnetic field interacts with the Earth's magnetic field and transfers an increased energy into the magnetosphere. Both interactions cause an increase in plasma movement through the magnetosphere (driven by increased electric fields inside the magnetosphere) and an increase in electric current in the magnetosphere and ionosphere. During the main phase of a geomagnetic storm, electric current in the magnetosphere creates a magnetic force that pushes out the boundary between the magnetosphere and the solar wind.
Several space weather phenomena tend to be associated with geomagnetic storms. These include solar energetic particle (SEP) events, geomagnetically induced currents (GIC), ionospheric storms and disturbances that cause radio and radar scintillation, disruption of navigation by magnetic compass and auroral displays at much lower magnetic latitudes than normal.
The largest recorded geomagnetic storm, the Carrington Event in September 1859, took down parts of the recently created US telegraph network, starting fires and electrically shocking telegraph operators. In 1989, a geomagnetic storm energized ground induced currents that disrupted electric power distribution throughout most of Quebec and caused aurorae as far south as Texas.
Definition
A geomagnetic storm is defined by changes in the Dst (disturbance – storm time) index. The Dst index estimates the globally averaged change of the horizontal component of the Earth's magnetic field at the magnetic equator based on measurements from a few magnetometer stations. Dst is computed once per hour and reported in near-real-time. During quiet times, Dst is between +20 and −20 nano-Tesla (nT).
A geomagnetic storm has three phases: initial, main and recovery. The initial phase is characterized by Dst (or its one-minute component SYM-H) increasing by 20 to 50 nT in tens of minutes. The initial phase is also referred to as a storm sudden commencement (SSC). However, not all geomagnetic storms have an initial phase and not all sudden increases in Dst or SYM-H are followed by a geomagnetic storm. The main phase of a geomagnetic storm is defined by Dst decreasing to less than −50 nT. The selection of −50 nT to define a storm is somewhat arbitrary. The minimum value during a storm will be between −50 and approximately −600 nT. The duration of the main phase is typically 2–8 hours. The recovery phase is when Dst changes from its minimum value to its quiet time value. The recovery phase may last as short as 8 hours or as long as 7 days.
The size of a geomagnetic storm is classified as moderate (−50 nT > minimum of Dst > −100 nT), intense (−100 nT > minimum Dst > −250 nT) or super-storm (minimum of Dst < −250 nT).
Measuring intensity
Geomagnetic storm intensity is reported in several different ways, including:
K-index
A-index
The G-scale used by the U.S. National Oceanic and Atmospheric Administration, which rates the storm from G1 to G5 (i.e. G1, G2, G3, G4, G5 in order), where G1 is the weakest storm classification (corresponding to a Kp value of 5), and G5 is the strongest (corresponding to a Kp value of 9).
History of the theory
In 1930, Sydney Chapman and Vincenzo C. A. Ferraro wrote an article, A New Theory of Magnetic Storms, that sought to explain the phenomenon. They argued that whenever the Sun emits a solar flare it also emits a plasma cloud, now known as a coronal mass ejection. They postulated that this plasma travels at a velocity such that it reaches Earth within 113 days, though we now know this journey takes 1 to 5 days. They wrote that the cloud then compresses the Earth's magnetic field and thus increases this field at the Earth's surface. Chapman and Ferraro's work drew on that of, among others, Kristian Birkeland, who had used recently-discovered cathode-ray tubes to show that the rays were deflected towards the poles of a magnetic sphere. He theorised that a similar phenomenon was responsible for auroras, explaining why they are more frequent in polar regions.
Occurrences
The first scientific observation of the effects of a geomagnetic storm occurred early in the 19th century: from May 1806 until June 1807, Alexander von Humboldt recorded the bearing of a magnetic compass in Berlin. On 21 December 1806, he noticed that his compass had become erratic during a bright auroral event.
On September 1–2, 1859, the largest recorded geomagnetic storm occurred. From August 28 until September 2, 1859, numerous sunspots and solar flares were observed on the Sun, with the largest flare on September 1. This is referred to as the solar storm of 1859 or the Carrington Event. It can be assumed that a massive coronal mass ejection was launched from the Sun and reached the Earth within eighteen hours—a trip that normally takes three to four days. The horizontal field was reduced by 1600 nT as recorded by the Colaba Observatory. It is estimated that Dst would have been approximately −1760 nT. Telegraph wires in both the United States and Europe experienced induced voltage increases (emf), in some cases even delivering shocks to telegraph operators and igniting fires. Aurorae were seen as far south as Hawaii, Mexico, Cuba and Italy—phenomena that are usually only visible in polar regions. Ice cores show evidence that events of similar intensity recur at an average rate of approximately once per 500 years.
Since 1859, less severe storms have occurred, notably the aurora of November 17, 1882 and the May 1921 geomagnetic storm, both with disruption of telegraph service and initiation of fires, and 1960, when widespread radio disruption was reported.
In early August 1972, a series of flares and solar storms peaks with a flare estimated around X20 producing the fastest CME transit ever recorded and a severe geomagnetic and proton storm that disrupted terrestrial electrical and communications networks, as well as satellites (at least one made permanently inoperative), and spontaneously detonated numerous U.S. Navy magnetic-influence sea mines in North Vietnam.
The March 1989 geomagnetic storm caused the collapse of the Hydro-Québec power grid in seconds as equipment protection relays tripped in a cascading sequence. Six million people were left without power for nine hours. The storm caused auroras as far south as Texas and Florida. The storm causing this event was the result of a coronal mass ejected from the Sun on March 9, 1989. The minimum Dst was −589 nT.
On July 14, 2000, an X5 class flare erupted (known as the Bastille Day event) and a coronal mass was launched directly at the Earth. A geomagnetic super storm occurred on July 15–17; the minimum of the Dst index was −301 nT. Despite the storm's strength, no power distribution failures were reported. The Bastille Day event was observed by Voyager 1 and Voyager 2, thus it is the farthest out in the Solar System that a solar storm has been observed.
Seventeen major flares erupted on the Sun between 19 October and 5 November 2003, including perhaps the most intense flare ever measured on the GOES XRS sensor—a huge X28 flare, resulting in an extreme radio blackout, on 4 November. These flares were associated with CME events that caused three geomagnetic storms between 29 October and 2 November, during which the second and third storms were initiated before the previous storm period had fully recovered. The minimum Dst values were −151, −353 and −383 nT. Another storm in this sequence occurred on 4–5 November with a minimum Dst of −69 nT. The last geomagnetic storm was weaker than the preceding storms, because the active region on the Sun had rotated beyond the meridian where the central portion CME created during the flare event passed to the side of the Earth. The whole sequence became known as the Halloween Solar Storm. The Wide Area Augmentation System (WAAS) operated by the Federal Aviation Administration (FAA) was offline for approximately 30 hours due to the storm. The Japanese ADEOS-2 satellite was severely damaged and the operation of many other satellites were interrupted due to the storm.
Interactions with planetary processes
The solar wind also carries with it the Sun's magnetic field. This field will have either a North or South orientation. If the solar wind has energetic bursts, contracting and expanding the magnetosphere, or if the solar wind takes a southward polarization, geomagnetic storms can be expected. The southward field causes magnetic reconnection of the dayside magnetopause, rapidly injecting magnetic and particle energy into the Earth's magnetosphere.
During a geomagnetic storm, the ionosphere's F2 layer becomes unstable, fragments, and may even disappear. In the northern and southern pole regions of the Earth, auroras are observable.
Instruments
Magnetometers monitor the auroral zone as well as the equatorial region. Two types of radar, coherent scatter and incoherent scatter, are used to probe the auroral ionosphere. By bouncing signals off ionospheric irregularities, which move with the field lines, one can trace their motion and infer magnetospheric convection.
Spacecraft instruments include:
Magnetometers, usually of the flux gate type. Usually these are at the end of booms, to keep them away from magnetic interference by the spacecraft and its electric circuits.
Electric sensors at the ends of opposing booms are used to measure potential differences between separated points, to derive electric fields associated with convection. The method works best at high plasma densities in low Earth orbit; far from Earth long booms are needed, to avoid shielding-out of electric forces.
Radio sounders from the ground can bounce radio waves of varying frequency off the ionosphere, and by timing their return determine the electron density profile—up to its peak, past which radio waves no longer return. Radio sounders in low Earth orbit aboard the Canadian Alouette 1 (1962) and Alouette 2 (1965), beamed radio waves earthward and observed the electron density profile of the "topside ionosphere". Other radio sounding methods were also tried in the ionosphere (e.g. on IMAGE).
Particle detectors include a Geiger counter, as was used for the original observations of the Van Allen radiation belt. Scintillator detectors came later, and still later "channeltron" electron multipliers found particularly wide use. To derive charge and mass composition, as well as energies, a variety of mass spectrograph designs were used. For energies up to about 50 keV (which constitute most of the magnetospheric plasma) time-of-flight spectrometers (e.g. "top-hat" design) are widely used.
Computers have made it possible to bring together decades of isolated magnetic observations and extract average patterns of electrical currents and average responses to interplanetary variations. They also run simulations of the global magnetosphere and its responses, by solving the equations of magnetohydrodynamics (MHD) on a numerical grid. Appropriate extensions must be added to cover the inner magnetosphere, where magnetic drifts and ionospheric conduction need to be taken into account. At polar regions, directly linked to the solar wind, large-scale ionospheric anomalies can be successfully modeled, even during geomagnetic super-storms. At smaller scales (comparable to a degree of latitude/longitude) the results are difficult to interpret, and certain assumptions about the high-latitude forcing uncertainty are needed.
Geomagnetic storm effects
Disruption of electrical systems
It has been suggested that a geomagnetic storm on the scale of the solar storm of 1859 today would cause billions or even trillions of dollars of damage to satellites, power grids and radio communications, and could cause electrical blackouts on a massive scale that might not be repaired for weeks, months, or even years. Such sudden electrical blackouts may threaten food production.
Main electrical grid
When magnetic fields move about in the vicinity of a conductor such as a wire, a geomagnetically induced current is produced in the conductor. This happens on a grand scale during geomagnetic storms (the same mechanism also influenced telephone and telegraph lines before fiber optics, see above) on all long transmission lines. Long transmission lines (many kilometers in length) are thus subject to damage by this effect. Notably, this chiefly includes operators in China, North America, and Australia, especially in modern high-voltage, low-resistance lines. The European grid consists mainly of shorter transmission circuits, which are less vulnerable to damage.
The (nearly direct) currents induced in these lines from geomagnetic storms are harmful to electrical transmission equipment, especially transformers—inducing core saturation, constraining their performance (as well as tripping various safety devices), and causing coils and cores to heat up. In extreme cases, this heat can disable or destroy them, even inducing a chain reaction that can overload transformers. Most generators are connected to the grid via transformers, isolating them from the induced currents on the grid, making them much less susceptible to damage due to geomagnetically induced current. However, a transformer that is subjected to this will act as an unbalanced load to the generator, causing negative sequence current in the stator and consequently rotor heating.
A 2008 study by Metatech corporation concluded that a storm with a strength comparable to that of 1921 would destroy more than 300 transformers and leave over 130 million people without power in the United States, costing several trillion dollars. The extent of the disruption is debated, with some congressional testimony indicating a potentially indefinite outage until transformers can be replaced or repaired. These predictions are contradicted by a North American Electric Reliability Corporation report that concludes that a geomagnetic storm would cause temporary grid instability but no widespread destruction of high-voltage transformers. The report points out that the widely quoted Quebec grid collapse was not caused by overheating transformers but by the near-simultaneous tripping of seven relays. In 2016, the United States Federal Energy Regulatory Commission adopted NEARC rules for equipment testing for electric utilities. Implementation of any upgrades needed to protect against the effects of geomagnetic storms was required within four years, and the regulations also directed further research.
Besides the transformers being vulnerable to the effects of a geomagnetic storm, electricity companies can also be affected indirectly by the geomagnetic storm. For instance, Internet service providers may go down during geomagnetic storms (and/or remain non-operational long after). Electricity companies may have equipment requiring a working Internet connection to function, so during the period the Internet service provider is down, the electricity too may not be distributed.
By receiving geomagnetic storm alerts and warnings (e.g. by the Space Weather Prediction Center; via Space Weather satellites as SOHO or ACE), power companies can minimize damage to power transmission equipment, by momentarily disconnecting transformers or by inducing temporary blackouts. Preventive measures also exist, including preventing the inflow of GICs into the grid through the neutral-to-ground connection.
Communications
High frequency (3–30 MHz) communication systems use the ionosphere to reflect radio signals over long distances. Ionospheric storms can affect radio communication at all latitudes. Some frequencies are absorbed and others are reflected, leading to rapidly fluctuating signals and unexpected propagation paths. TV and commercial radio stations are little affected by solar activity, but ground-to-air, ship-to-shore, shortwave broadcast and amateur radio (mostly the bands below 30 MHz) are frequently disrupted. Radio operators using HF bands rely upon solar and geomagnetic alerts to keep their communication circuits up and running.
Military detection or early warning systems operating in the high frequency range are also affected by solar activity. The over-the-horizon radar bounces signals off the ionosphere to monitor the launch of aircraft and missiles from long distances. During geomagnetic storms, this system can be severely hampered by radio clutter. Also some submarine detection systems use the magnetic signatures of submarines as one input to their locating schemes. Geomagnetic storms can mask and distort these signals.
The Federal Aviation Administration routinely receives alerts of solar radio bursts so that they can recognize communication problems and avoid unnecessary maintenance. When an aircraft and a ground station are aligned with the Sun, high levels of noise can occur on air-control radio frequencies. This can also happen on UHF and SHF satellite communications, when an Earth station, a satellite and the Sun are in alignment. In order to prevent unnecessary maintenance on satellite communications systems aboard aircraft AirSatOne provides a live feed for geophysical events from NOAA's Space Weather Prediction Center. allows users to view observed and predicted space storms. Geophysical Alerts are important to flight crews and maintenance personnel to determine if any upcoming activity or history has or will have an effect on satellite communications, GPS navigation and HF Communications.
Telegraph lines in the past were affected by geomagnetic storms. Telegraphs used a single long wire for the data line, stretching for many miles, using the ground as the return wire and fed with DC power from a battery; this made them (together with the power lines mentioned below) susceptible to being influenced by the fluctuations caused by the ring current. The voltage/current induced by the geomagnetic storm could have diminished the signal, when subtracted from the battery polarity, or to overly strong and spurious signals when added to it; some operators learned to disconnect the battery and rely on the induced current as their power source. In extreme cases the induced current was so high the coils at the receiving side burst in flames, or the operators received electric shocks. Geomagnetic storms affect also long-haul telephone lines, including undersea cables unless they are fiber optic.
Damage to communications satellites can disrupt non-terrestrial telephone, television, radio and Internet links. The National Academy of Sciences reported in 2008 on possible scenarios of widespread disruption in the 2012–2013 solar peak. A solar superstorm could cause large-scale global months-long Internet outages. A study describes potential mitigation measures and exceptions – such as user-powered mesh networks, related peer-to-peer applications and new protocols – and analyzes the robustness of the current Internet infrastructure.
Navigation systems
Global navigation satellite systems (GNSS), and other navigation systems such as LORAN and the now-defunct OMEGA are adversely affected when solar activity disrupts their signal propagation. The OMEGA system consisted of eight transmitters located throughout the world. Airplanes and ships used the very low frequency signals from these transmitters to determine their positions. During solar events and geomagnetic storms, the system gave navigators information that was inaccurate by as much as several miles. If navigators had been alerted that a proton event or geomagnetic storm was in progress, they could have switched to a backup system.
GNSS signals are affected when solar activity causes sudden variations in the density of the ionosphere, causing the satellite signals to scintillate (like a twinkling star). The scintillation of satellite signals during ionospheric disturbances is studied at HAARP during ionospheric modification experiments. It has also been studied at the Jicamarca Radio Observatory.
One technology used to allow GNSS receivers to continue to operate in the presence of some confusing signals is Receiver Autonomous Integrity Monitoring (RAIM), used by GPS. However, RAIM is predicated on the assumption that a majority of the GPS constellation is operating properly, and so it is much less useful when the entire constellation is perturbed by global influences such as geomagnetic storms. Even if RAIM detects a loss of integrity in these cases, it may not be able to provide a useful, reliable signal.
Satellite hardware damage
Geomagnetic storms and increased solar ultraviolet emission heat Earth's upper atmosphere, causing it to expand. The heated air rises, and the density at the orbit of satellites up to about increases significantly. This results in increased drag, causing satellites to slow and change orbit slightly. Low Earth orbit satellites that are not repeatedly boosted to higher orbits slowly fall and eventually burn up. Skylab's 1979 destruction is an example of a spacecraft reentering Earth's atmosphere prematurely as a result of higher-than-expected solar activity. During the great geomagnetic storm of March 1989, four of the U.S. Navy's navigational satellites had to be taken out of service for up to a week, the U.S. Space Command had to post new orbital elements for over 1000 objects affected, and the Solar Maximum Mission satellite fell out of orbit in December the same year.
The vulnerability of the satellites depends on their position as well. The South Atlantic Anomaly is a perilous place for a satellite to pass through, due to the unusually weak geomagnetic field at low Earth orbit.
Pipelines
Rapidly fluctuating geomagnetic fields can produce geomagnetically induced currents in pipelines. This can cause multiple problems for pipeline engineers. Pipeline flow meters can transmit erroneous flow information and the corrosion rate of the pipeline can be dramatically increased.
Radiation hazards to humans
Earth's atmosphere and magnetosphere allow adequate protection at ground level, but astronauts are subject to potentially lethal radiation poisoning. The penetration of high-energy particles into living cells can cause chromosome damage, cancer and other health problems. Large doses can be immediately fatal. Solar protons with energies greater than 30 MeV are particularly hazardous.
Solar proton events can also produce elevated radiation aboard aircraft flying at high altitudes. Although these risks are small, flight crews may be exposed repeatedly, and monitoring of solar proton events by satellite instrumentation allows exposure to be monitored and evaluated, and eventually flight paths and altitudes to be adjusted to lower the absorbed dose.
Ground level enhancements, also known as ground level events or GLEs, occur when a solar particle event contains particles with sufficient energy to have effects at ground level, mainly detected as an increase in the number of neutrons measured at ground level. These events have been shown to have an impact on radiation dosage, but they do not significantly increase the risk of cancer.
Effect on animals
There is a large but controversial body of scientific literature on connections between geomagnetic storms and human health. This began with Russian papers, and the subject was subsequently studied by Western scientists. Theories for the cause include the involvement of cryptochrome, melatonin, the pineal gland, and the circadian rhythm.
Some scientists suggest that solar storms induce whales to beach themselves. Some have speculated that migrating animals which use magnetoreception to navigate, such as birds and honey bees, might also be affected.
| Physical sciences | Storms | Earth science |
147699 | https://en.wikipedia.org/wiki/Laundry | Laundry | Laundry is the washing of clothing and other textiles, and, more broadly, their drying and ironing as well. Laundry has been part of history since humans began to wear clothes, so the methods by which different cultures have dealt with this universal human need are of interest to several branches of scholarship.
Laundry work has traditionally been highly gendered, with the responsibility in most cultures falling to women (formerly known as laundresses or washerwomen). The Industrial Revolution gradually led to mechanized solutions to laundry work, notably the washing machine and later the tumble dryer. Laundry, like cooking and child care, is still done both at home and by commercial establishments outside the home.
The word "laundry" may refer to the clothing itself, or to the place where the cleaning happens. An individual home may have a laundry room; a utility room includes, but is not restricted to, the function of washing clothes. An apartment building or student hall of residence may have a shared laundry facility such as a tvättstuga. A stand-alone business is referred to as a self-service laundry (launderette in British English or laundromat in North American English).
History
Watercourses
Laundry was first done in watercourses, letting the water carry away the materials which could cause stains and smells. Laundry is still done this way in the rural regions of poor countries. Agitation helps remove the dirt, so the laundry was rubbed, twisted, or slapped against flat rocks. One name for this surface is a beetling-stone, related to beetling, a technique in the production of linen; one name for a wooden substitute is a battling-block. The dirt was beaten out with a wooden implement known as a washing paddle, battling stick, bat, beetle or club. Wooden or stone scrubbing surfaces set up near a water supply were gradually replaced by portable rub boards, eventually factory-made corrugated glass or metal washboards.
Once clean, the clothes were wrung out — twisted to remove most of the water. Then they were hung up on poles or clothes lines to air dry, or sometimes just spread out on clean grass, bushes, or trees.
Washhouses
Before the advent of the washing machine, laundry was often done in a communal setting.
Villages across Europe that could afford it built a wash-house, sometimes known by the French name of lavoir. Water was channelled from a stream or spring and fed into a building, possibly just a roof with no walls. This wash-house usually contained two basins – one for washing and the other for rinsing – through which the water was constantly flowing, as well as a stone lip inclined towards the water against which the wet laundry could be beaten. Such facilities were more comfortable and convenient than washing in a watercourse. Some lavoirs had the wash-basins at waist height, although others remained on the ground. The launderers were protected to some extent from rain, and their travel was reduced, as the facilities were usually at hand in the village or at the edge of a town. These facilities were public and available to all families, and usually used by the entire village. Many of these village wash-houses are still standing, historic structures with no obvious modern purpose.
The job of doing the laundry was reserved for women, who washed all their family's laundry. Washerwomen (laundresses) took in the laundry of others, charging by the piece. As such, wash-houses were an obligatory stop in many women's weekly lives and became a sort of institution or meeting place. It was a women-only space where they could discuss issues or simply chat (cf the concept of the village pump). Indeed, this tradition is reflected in the Catalan idiom "fer safareig" (literally, "to do the laundry"), which means to gossip.
European cities also had public wash-houses. The city authorities wanted to give the poorer population, who would otherwise not have access to laundry facilities, the opportunity to wash their clothes. Sometimes these facilities were combined with public baths, see for example Baths and wash houses in Britain. The aim was to foster hygiene and thus reduce outbreaks of epidemics.
Sometimes large metal cauldrons (a "wash copper", even when not made of that metal), were filled with fresh water and heated over a fire, as hot or boiling water is more effective than cold in removing dirt. A posser could be used to agitate clothes in a tub. A related implement called a washing dolly is "a wooden stick or mallet with an attached cluster of legs or pegs" that moves the cloth through the water.
Washing machines and other devices
The Industrial Revolution completely transformed laundry technology. Christina Hardyment, in her history from the Great Exhibition of 1851, argues that it was the development of domestic machinery that led to women's liberation.
The mangle (or "wringer" in American English) was developed in the 19th century — two long rollers in a frame and a crank to revolve them. A laundry-worker took sopping wet clothing and cranked it through the mangle, compressing the cloth and expelling the excess water. The mangle was much quicker than hand twisting. It was a variation on the box mangle used primarily for pressing and smoothing cloth.
Meanwhile, 19th-century inventors further mechanized the laundry process with various hand-operated washing machines to replace tedious hand rubbing against a washboard. Most involved turning a handle to move paddles inside a tub. Then some early-20th-century machines used an electrically powered agitator. Many of these washing machines were simply a tub on legs, with a hand-operated mangle on top. Later the mangle too was electrically powered, then replaced by a perforated double tub, which spun out the excess water in a spin cycle.
Laundry drying was also mechanized, with clothes dryers. Dryers were also spinning perforated tubs, but they blew heated air rather than water.
Chinese laundries in North America
In the late 19th and early 20th century, Chinese immigrants to the United States and to Canada were well represented as laundry workers. Discrimination, lack of English-language skills, and lack of capital kept Chinese immigrants out of most desirable careers. Around 1900, one in four ethnic Chinese men in the U.S. worked in a laundry, typically working 10 to 16 hours a day.
Chinese people in New York City were running an estimated 3,550 laundries at the beginning of the Great Depression. In 1933, the city's Board of Aldermen passed a law clearly intended to drive the Chinese out of the business. Among other things, it limited ownership of laundries to U.S. citizens. The Chinese Consolidated Benevolent Association tried fruitlessly to fend this off, resulting in the formation of the openly leftist Chinese Hand Laundry Alliance (CHLA), which successfully challenged this provision of the law, allowing Chinese laundry workers to preserve their livelihoods. The CHLA went on to function as a more general civil rights group; its numbers declined strongly after it was targeted by the FBI during the Second Red Scare (1947–1957).
South Africa
From 1850 to 1910, Zulu men took on the task of laundering the clothes of Europeans, both Boers and British. "Laundering recalled the specialist craft of hide-dressing in which Zulu males engaged as izinyanga, a prestige occupation that paid handsomely." They created a guild structure, similar to a union, to guard their conditions and wages, evolving into "one, if not indeed the most, powerful group of African work-men in nineteenth-century Natal".
India
In India, laundry was traditionally done by men. A washerman was called a dhobiwallah, and dhobi became the name of their caste group.
A laundry-place is generally called a dhobi ghat; this has given rise to place names where they work or worked, including Mahalaxmi Dhobi Ghat in Mumbai, Dhoby Ghaut in Singapore and Dhobi Ghaut in Penang, Malaysia.
Philippines
Until the early 1980s, when washing machines became more affordable in the country, much of the laundry work in the Philippines was done manually, and this role was generally assigned to women. A professional laundrywoman was called a labandera.
Ancient Rome
The workers in ancient Rome who cleaned the cloth were called fullones, singular fullo (cf fulling, a process in wool-making, and Fuller's earth, used to clean). Clothes were treated in small tubs standing in niches surrounded by low walls, known as treading or fulling stalls. The tub was filled with water and a mixture of alkaline chemicals (sometimes including urine). The fuller stood in the tub and trampled the cloth, a technique known elsewhere as posting. The aim of this treatment was to apply the chemical agents to the cloth so that they could do their work, the resolving of greases and fats. These stalls are so typical of these workshops that they are used to identify fullonicae in the archaeological remains.
Laundry processes
Laundry processes include washing (usually with water containing detergents or other chemicals), agitation, rinsing, drying, pressing (ironing), and folding. The washing will sometimes be done at a temperature above room temperature to increase the activities of any chemicals used and the solubility of stains, and high temperatures kill micro-organisms that may be present on the fabric. However, it is advised that cotton be washed at a cooler temperature to prevent shrinking. Many professional laundry services are present in the market which offers at different price range.
Agitation helps remove dirt which is usually mobilised by surfactants from between fibres, however, due to the small size of the pores in fibres, the 'stagnant core' of the fibres themselves see virtually no flow. The fibres are nevertheless rapidly cleaned by diffusiophoresis carrying dirt out into the clean water during the rinsing process.
Chemicals
Various chemicals may be used to increase the solvent power of water, such as the compounds in soaproot or yucca-root used by Native American tribes, or the ash lye (usually sodium hydroxide or potassium hydroxide) once widely used for soaking laundry in Europe. Soap, a compound made from lye and fat, is an ancient and common laundry aid. Modern washing machines typically use synthetic powdered or liquid laundry detergent in place of more traditional soap.
Cleaning or dry cleaning
Dry cleaning refers to any process which uses a chemical solvent other than water. The solvent used is typically tetrachloroethylene (perchloroethylene), which the industry calls "perc". It is used to clean delicate fabrics that cannot withstand the rough and tumble of a washing machine and clothes dryer; it can also obviate labor-intensive hand washing.
Shared laundry rooms
In some parts of the world, including North America, apartment buildings and dormitories often have laundry rooms, where residents share washing machines and dryers. Usually the machines are set to run only when money is put in a coin slot.
In other parts of the world, including Europe, apartment buildings with laundry rooms are uncommon, and each apartment may have its own washing machine. Those without a machine at home or the use of a laundry room must either wash their clothes by hand or visit a commercial self-service laundry (laundromat, laundrette) or a laundry shop, such as 5àsec.
Right to dry movement
Some American communities forbid their residents from drying clothes outside, and citizens protesting this have created a "right to dry" movement. Many homeowners' associations and other communities in the United States prohibit residents from using a clothesline outdoors, or limit such use to locations that are not visible from the street or to certain times of day. Other communities, however, expressly prohibit rules that prevent the use of clotheslines. Some organizations have been campaigning against legislation which has outlawed line-drying of clothing in public places, especially given the increased greenhouse gas emissions produced by some types of electrical power generation needed to power electric clothes dryers, since driers can constitute a considerable fraction of a home's total energy usage.
Florida ("the Sunshine State") is the only state to expressly guarantee a right to dry, although Utah and Hawaii have passed solar rights legislation. A Florida law explicitly states: "No deed restrictions, covenants, or similar binding agreements running with the land shall prohibit or have the effect of prohibiting solar collectors, clotheslines, or other energy devices based on renewable resources from being installed on buildings erected on the lots or parcels covered by the deed restrictions, covenants, or binding agreements." No other state has such clearcut legislation. Vermont considered a "Right to Dry" bill in 1999, but it was defeated in the Senate Natural Resources & Energy Committee. The language has been included in a 2007 voluntary energy conservation bill, introduced by Senator Dick McCormack. Legislation making it possible for thousands of American families to start using clotheslines in communities where they were formerly banned was passed in Colorado in 2008. In 2009, clothesline legislation was debated in the states of Connecticut, Hawaii, Maryland, Maine, New Hampshire, Nebraska, Oregon, Virginia, and Vermont.
Similar measures have been introduced in Canada - in particular, the province of Ontario.
Common problems
Novice users of modern laundry machines sometimes experience accidental shrinkage of garments, especially when applying heat. For wool garments, this is due to scales on the fibers, which heat and agitation cause to stick together. Other fabrics (like cotton) have their fibers stretched by mechanical force during production, and can shrink slightly when heated (though to a lesser degree than wool). Some clothes are "pre-shrunk" to avoid this problem.
Another common problem is color bleeding from dyed articles to white or pale-colored ones. Many laundry guides suggest washing whites separately from colored items. Sometimes only similar colors are washed together to avoid this problem, which is lessened by cold water and repeated washings. Sometimes this blending of colors is seen as a selling point, as with madras cloth.
Laundry symbols are included on many clothes to help consumers avoid these problems.
Synthetic fibers in laundry can also contribute to microplastic pollution.
Etymology
The word laundry comes from Middle English lavendrye, laundry, from Old French lavanderie, from lavandier.
In culture
In Homer's Odyssey, Princess Nausicaa and her handmaidens are washing laundry by the shore when they see and rescue the ship-wrecked Ulysses.
| Technology | Food, water and health | null |
147728 | https://en.wikipedia.org/wiki/Toothbrush | Toothbrush | A toothbrush is a special type of brush used to clean the teeth, gums, and tongue. It consists of a head of tightly clustered bristles, atop of which toothpaste can be applied, mounted on a handle which facilitates the cleaning of hard-to-reach areas of the mouth. They should be used in conjunction with something to clean between the teeth where the bristles of the toothbrush cannot reach - for example floss, tape or interdental brushes.
They are available with different bristle textures, sizes, and forms. Most dentists recommend using a soft toothbrush since hard-bristled toothbrushes can damage tooth enamel and irritate the gums.
Because many common and effective ingredients in toothpaste are harmful if swallowed in large doses, tooth paste should instead should be spat out. The act of brushing teeth is most often done at a sink within the kitchen or bathroom, where the brush may be rinsed off afterwards to remove any debris remaining and then dried to reduce conditions ideal for bacterial growth (and, if it is a wooden toothbrush, mold as well).
Some toothbrushes have plant-based handles, often bamboo. However, numerous others are made of cheap plastic; such brushes constitute a significant source of pollution. Over 1 billion toothbrushes are disposed of into landfills annually in the United States alone. Bristles are commonly made of nylon (which, while not biodegradable, as plastic is, may still be recycled), bamboo viscose, or bristle of boar.
History
Precursors
Before the invention of the toothbrush, a variety of oral hygiene measures had been used. This has been verified by excavations during which tree twigs, bird feathers, animal bones and porcupine quills were recovered.
The predecessor of the toothbrush is the chew stick. Chew sticks were twigs with frayed ends used to brush the teeth while the other end was used as a toothpick. The earliest chew sticks were discovered in Sumer in southern Mesopotamia in 3500 BC, an Egyptian tomb dating from 3000 BC, and mentioned in Chinese records dating from 1600 BC.
The Indian way of using tooth wood for brushing is presented by the Chinese Monk Yijing (635–713 CE) when he describes the rules for monks in his book:
"Every day in the morning, a monk must chew a piece of tooth wood to brush his teeth and scrape his tongue, and this must be done in the proper way. Only after one has washed one's hands and mouth may one make salutations. Otherwise both the saluter and the saluted are at fault. In Sanskrit, the tooth wood is known as the dantakastha—danta meaning tooth, and kastha, a piece of wood. It is twelve finger-widths in length. The shortest is not less than eight finger-widths long, resembling the little finger in size. Chew one end of the wood well for a long while and then brush the teeth with it."
The Greeks and Romans used toothpicks to clean their teeth, and toothpick-like twigs have been excavated in Qin dynasty tombs. Chew sticks remain common in Africa, the rural Southern United States, and in the Islamic world, the use of the chewing stick, miswak, is regarded as a pious action and is prescribed for use before every prayer, occurring five times a day. Miswaks have been used by Muslims since the 7th century. Twigs of Neem Tree have been used by ancient Indians. Neem, in its full bloom, can aid in healing by keeping the area clean and disinfected. In fact, even today, Neem twigs called datun are used for brushing teeth in India, although not hugely common.
Toothbrush
The first bristle toothbrush resembling the modern one was found in China. Used during the Tang dynasty (619–907), it consisted of hog bristles. The bristles were sourced from hogs living in Siberia and northern China because the colder temperatures provided firmer bristles. They were attached to a handle manufactured from bamboo or bone, forming a toothbrush. In 1223, Japanese Zen master Dōgen Kigen recorded in his Shōbōgenzō that he saw monks in China clean their teeth with brushes made of horsetail hairs attached to an oxbone handle. The bristle toothbrush spread to Europe, brought from China to Europe by travellers. It was adopted in Europe during the 17th century. The earliest identified use of the word toothbrush in English was in the autobiography of Anthony Wood who wrote in 1690 that he had bought a toothbrush from J. Barret. Europeans found the hog bristle toothbrushes imported from China too firm and preferred softer bristle toothbrushes made from horsehair. Mass-produced toothbrushes made with horse or boar bristle continued to be imported to Britain from China until the mid 20th century.
In the UK, William Addis is believed to have produced the first mass-produced toothbrush in 1780. In 1770, he had been jailed for causing a riot. While in prison he decided that using a rag with soot and salt on the teeth was ineffective and could be improved. After saving a small bone from a meal, he drilled small holes into the bone and tied into the bone tufts of bristles that he had obtained from one of the guards, passed the tufts of bristle through the holes in the bone and sealed the holes with glue. After his release, he became wealthy after starting a business manufacturing toothbrushes. He died in 1808, bequeathing the business to his eldest son. It remained within family ownership until 1996. Under the name Wisdom Toothbrushes, the company now manufactures 70 million toothbrushes per year in the UK. By 1840 toothbrushes were being mass-produced in Britain, France, Germany, and Japan. Pig bristles were used for cheaper toothbrushes and badger hair for the more expensive ones.
Hertford Museum in Hertford, UK, holds approximately 5000 brushes that make up part of the Addis Collection. The Addis factory on Ware Road was a major employer in the town until 1996. Since the closure of the factory, Hertford Museum has received photographs and documents relating to the archive, and collected oral histories from former employees.
The first patent for a toothbrush was granted to H.N. Wadsworth in 1857 (U.S.A. Patent No. 18,653) in the United States, but mass production in the United States did not start until 1885. The improved design had a bone handle with holes bored into it for the Siberian boar hair bristles. Unfortunately, animal bristle was not an ideal material as it retained bacteria, did not dry efficiently and the bristles often fell out. In addition to bone, handles were made of wood or ivory. In the United States, brushing teeth did not become routine until after World War II, when American soldiers had to clean their teeth daily.
During the 1900s, celluloid gradually replaced bone handles. Natural animal bristles were also replaced by synthetic fibers, usually nylon, by DuPont in 1938. The first nylon bristle toothbrush made with nylon yarn went on sale on February 24, 1938. The first electric toothbrush, the Broxodent, was invented in Switzerland in 1954. By the turn of the 21st century nylon had come to be widely used for the bristles and the handles were usually molded from thermoplastic materials.
Johnson & Johnson, a leading medical supplies firm, introduced the "Reach" toothbrush in 1977. It differed from previous toothbrushes in three ways: it had an angled head, similar to dental instruments, to reach back teeth; the bristles were concentrated more closely than usual to clean each tooth of potentially cariogenic (cavity-causing) materials; and the outer bristles were longer and softer than the inner bristles. Other manufacturers soon followed with other designs aimed at improving effectiveness. In spite of the changes with the number of tufts and the spacing, the handle form and design, the bristles were still straight and difficult to maneuver. In 1978 Dr. George C. Collis developed the Collis Curve toothbrush which was the first toothbrush to have curved bristles. The curved bristles follow the curvature of the teeth and safely reach in between the teeth and into the sulcular areas.
In January 2003, the toothbrush was selected as the number one invention Americans could not live without according to the Lemelson-MIT Invention Index.
Types of toothbrush
Multi-sided toothbrushes
A multi-sided toothbrush is a fast and easy way to brush the teeth.
Electric toothbrush
It has been discovered that compared to a manual brush, the multi-directional power brush might reduce the incidence of gingivitis and plaque, when compared to regular side-to-side brushing. These brushes tend to be more costly and damaging to the environment when compared to manual toothbrushes. Most studies report performances equivalent to those of manual brushings, possibly with a decrease in plaque and gingivitis. An additional timer and pressure sensors can encourage a more efficient cleaning process.
Electric toothbrushes can be classified, according to the speed of their movements as: standard power toothbrushes, sonic toothbrushes, or ultrasonic toothbrushes. Any electric toothbrush is technically a powered toothbrush. If the motion of the toothbrush is sufficiently rapid to produce a hum in the audible frequency range (20 Hz to 20,000 Hz), it can be classified as a sonic toothbrush. Any electric toothbrush with movement faster than this limit can be classified as an ultrasonic toothbrush. Certain ultrasonic toothbrushes, such as the Megasonex and the Ultreo, have both sonic and ultrasonic movements.
There are different electric toothbrush heads designed for sensitive teeth and gums, increased stain removal, or different-sized bristles for tight or gapped teeth. The hand motion with an electric toothbrush is different from a manual toothbrush. They are meant to have the bristles do the work by just placing and moving the toothbrush, so that fewer back and forth strokes are needed.
Interdental brush
An interdental or interproximal ("proxy") brush is a small brush, typically disposable, either supplied with a reusable angled plastic handle or an integral handle, used for cleaning between teeth and between the wires of dental braces and the teeth.
The use of interdental brushes in conjunction with tooth brushing has been shown to reduce both the amount of plaque and the incidence of gingivitis when compared to tooth brushing alone. Although there is some evidence that after tooth brushing with a conventional tooth brush, interdental brushes remove more plaque than dental floss, a systematic review reported insufficient evidence to determine such an association.
The size of interdental brushes is standardized in ISO 16409. The brush size, which is a number between 0 (small space between teeth) and 8 (large space), indicates the passage hole diameter. This corresponds to the space between two teeth that is just sufficient for the brush to go through without bending the wire. The color of the brushes differs between producers. The same is the case with respect to the wire diameter.
End-tuft brush
The small round brush head comprises seven tufts of tightly packed soft nylon bristles, trimmed so the bristles in the center can reach deeper into small spaces. The brush handle is ergonomically designed for a firm grip, giving the control and precision necessary to clean where most other cleaning aids cannot reach. These areas include the posterior of the wisdom teeth (third molars), orthodontic structures (braces), crowded teeth, and tooth surfaces that are next to missing teeth. It can also be used to clean areas around implants, bridges, dentures and other appliances.
Chewable toothbrush
A chewable toothbrush is a miniature plastic moulded toothbrush which can be placed inside the mouth. While not commonly used, they are useful to travelers and are sometimes available from bathroom vending machines. They are available in different flavors such as mint or bubblegum and should be disposed of after use. Other types of disposable toothbrushes include those that contain a small breakable plastic ball of toothpaste on the bristles, which can be used without water.
Musical toothbrush
A musical toothbrush is a type of manual or powered toothbrush designed to make tooth brushing habit more interesting. It is more commonly introduced to children to gain their attention and positively influence their tooth brushing behavior. The music starts while child starts brushing, it continuously plays during the brushing and it ends when the child stops brushing.
Toothpaste-less toothbrush
Reusable toothbrush
Tooth brushing
Hygiene and care
It is not recommended to share toothbrushes with others, since besides general hygienic concerns, there is a risk of transmitting diseases that are typically transmittable by blood, such as Hepatitis C.
It is advisable to rinse the toothbrush with water, shake it off and let dry after use.
Studies have shown that brushing to remove dental plaque more often than every 48 hours is enough to maintain gum and tooth health. Tooth brushing can remove plaque up to one millimeter below the gum line. Each person has a habitual brushing method, so more frequent brushing does not cover additional parts of the teeth or mouth.
Most dentists recommended patients brush twice a day in the hope that frequent brushing would clean more areas of the mouth.
Tooth brushing is the most common preventive healthcare activity, but tooth and gum disease remain high, since lay people clean at most 40% of their tooth margins at the gum line. Videos show that even when asked to brush their best, they do not know how to clean effectively.
Adversity of toothbrushes
Teeth can be damaged by several factors including poor oral hygiene, but also by wrong oral hygiene. Especially for sensitive teeth, damage to dentin and gums can be prevented by several measures including a correct brushing technique.
It is beneficial when using a straight bristled brush: not to scrub horizontally over the necks of teeth; not to press the brush too hard against the teeth, to choose a toothpaste that is not too abrasive; and to wait at least 30 minutes after consumption of acidic food or drinks before brushing.
Harder toothbrushes reduce plaque more efficiently but are more stressful to teeth and gum; using a medium to soft brush for a longer cleaning time was rated to be the best compromise between cleaning result and gum and tooth health.
A study by University College London found that advice on brushing technique and frequency given by 10 national dental associations, toothpaste and toothbrush companies, and in dental textbooks was inconsistent.
| Biology and health sciences | Hygiene products | Health |
147735 | https://en.wikipedia.org/wiki/Toothpaste | Toothpaste | Toothpaste is a paste or gel dentifrice used with a toothbrush to clean and maintain the aesthetics and health of teeth. Toothpaste is used to promote oral hygiene: it is an abrasive that aids in removing dental plaque and food from the teeth, assists in suppressing halitosis, and delivers active ingredients (most commonly fluoride) to help prevent tooth decay (dental caries) and gum disease (gingivitis). Owing to differences in composition and fluoride content, not all toothpastes are equally effective in maintaining oral health. The decline of tooth decay during the 20th century has been attributed to the introduction and regular use of fluoride-containing toothpastes worldwide. Large amounts of swallowed toothpaste can be poisonous. Common colors for toothpaste include white (sometimes with colored stripes or green tint) and blue.
History
Early toothpastes
Since 5000 BCE, the Egyptians made a tooth powder, which consisted of powdered ashes of ox hooves, myrrh, powdered and burnt eggshells, and pumice. The Greeks, and then the Romans, improved the recipes by adding abrasives such as crushed bones and oyster shells. In the 9th century, Iraqi musician and fashion designer Ziryab invented a type of toothpaste, which he popularized throughout Islamic Spain. The exact ingredients of this toothpaste are unknown, but it was reported to have been both "functional and pleasant to taste". It is not known whether these early toothpastes were used alone, were to be rubbed onto the teeth with rags, or were to be used with early toothbrushes, such as neem-tree twigs and miswak. During Japan's Edo period, inventor Hiraga Gennai's Hika rakuyo (1769) contained advertisements for Sosekiko, a "toothpaste in a box." Toothpastes or powders came into general use in the 19th century.
Usefulness
Toothpastes are generally useful to maintain dental health. Toothpastes containing fluoride are effective at preventing tooth decay. Toothpastes may also help to control and remove plaque build-up, promoting healthy gums. A 2016 systematic review indicated that using toothpaste when brushing the teeth does not necessarily impact the level of plaque removal. However, the active ingredients in toothpastes are able to prevent dental diseases with regular use.
Ingredients
Toothpastes are derived from a variety of components, the three main ones being abrasives, fluoride, and detergent.
Abrasives
Abrasives constitute 8-20% of a typical toothpaste. These insoluble particles are designed to help remove plaque from the teeth. The removal of plaque inhibits the accumulation of tartar (calculus) helping to minimize the risk of gum disease. Representative abrasives include particles of aluminum hydroxide (Al(OH)3), calcium carbonate (CaCO3), magnesium carbonate (MgCO3), sodium bicarbonate, various calcium hydrogen phosphates, various silicas and zeolites, and hydroxyapatite (Ca5(PO4)3OH).
After the Microbead-Free Waters Act of 2015, the use of microbeads in toothpaste has been discontinued in the US, however since 2015 the industry has shifted toward instead using FDA-approved "rinse-off" metallized-plastic glitter as their primary abrasive agent. Some brands contain powdered white mica, which acts as a mild abrasive, and also adds a cosmetic glittery shimmer to the paste. The polishing of teeth removes stains from tooth surfaces, but has not been shown to improve dental health over and above the effects of the removal of plaque and Calculus.
Abrasives, like the dental polishing agents used in dentists' offices, also cause a small amount of enamel erosion which is termed "polishing" action. The abrasive effect of toothpaste is indicated by its RDA value. Toothpastes with RDA values above 250 are potentially damaging to the surfaces of teeth. The American National Standards Institute and American Dental Association considers toothpastes with an RDA below 250 to be safe and effective for a lifetime of use.
Fluorides
Fluoride in various forms is the most popular and effective active ingredient in toothpaste to prevent cavities. Fluoride is present in small amounts in plants, animals, and some natural water sources. The additional fluoride in toothpaste has beneficial effects on the formation of dental enamel and bones. Sodium fluoride (NaF) is the most common source of fluoride, but stannous fluoride (SnF2), and sodium monofluorophosphate (Na2PO3F) are also used. At similar fluoride concentrations, toothpastes containing stannous fluoride have been shown to be more effective than toothpastes containing sodium fluoride for reducing the incidence of dental caries and dental erosion, as well as reducing gingivitis. Some stannous fluoride-containing toothpastes also contain ingredients that allow for better stain and calculus removal. A systematic review revealed stabilised stannous fluoride-containing toothpastes had a positive effect on the reduction of plaque, gingivitis and staining, with a significant reduction in calculus and halitosis compared to other toothpastes. Furthermore, numerous clinical trials have shown gluconate chelated stannous fluoride toothpastes possess superior protection against dental erosion and dentine hypersensitivity compared to other fluoride-containing and fluoride-free toothpastes.
Much of the toothpaste sold in the United States has 1,000 to 1,100 parts per million fluoride. In European countries, such as the UK or Greece, the fluoride content is often higher; a sodium fluoride content of 0.312% w/w (1,450 ppm fluoride) or stannous fluoride content of 0.454% w/w (1,100 ppm fluoride) is common. All of these concentrations are likely to prevent tooth decay, according to a 2019 Cochrane review. Concentrations below 1,000 ppm are not likely to be preventive, and the preventive effect increases with concentration. Clinical trials support the use of high fluoride (5,000 ppm fluoride) dentifrices, for prevention of root caries in elderly adults by reducing the amount of plaque accumulated, decreasing the number of mutans streptococci and lactobacilli and possibly promoting calcium fluoride deposits to a higher degree than after the use of traditional fluoride containing dentifrices.
Surfactants
Many, although not all, toothpastes contain sodium lauryl sulfate (SLS) or related surfactants (detergents). SLS is found in many other personal care products as well, such as shampoo, and is mainly a foaming agent, which enables uniform distribution of toothpaste, improving its cleansing power.
Other components
Antibacterial agents
Triclosan, an antibacterial agent, is a common toothpaste ingredient in the United Kingdom. Triclosan or zinc chloride prevent gingivitis and, according to the American Dental Association, helps reduce tartar and bad breath. A 2006 review of clinical research concluded there was evidence for the effectiveness of 0.30% triclosan in reducing plaque and gingivitis. Another Cochrane review in 2013 has found that triclosan achieved a 22% reduction in plaque, and in gingivitis, a 48% reduction in bleeding gums. However, there was insufficient evidence to show a difference in fighting periodontitis and there was no evidence either of any harmful effects associated with the use of triclosan toothpastes for more than 3 years. The evidence relating to plaque and gingivitis was considered to be of moderate quality while for periodontitis was low quality. Recently, triclosan has been removed as an ingredient from well-known toothpaste formulations. This may be attributed to concerns about adverse effects associated with triclosan exposure. Triclosan use in cosmetics has been positively correlated with triclosan levels in human tissues, plasma and breast milk, and is considered to have potential neurotoxic effects. Long-term studies are needed to substantiate these concerns.
Chlorhexidine is another antimicrobial agent used in toothpastes; however, it is more commonly added in mouthwash products. Sodium laureth sulfate, a foaming agent, is a common toothpaste ingredient that also possesses some antimicrobial activities. There are also many commercial products available in the market containing different essential oils, herbal ingredients (e.g. chamomile, neem, chitosan, Aloe vera), and natural or plant extracts (e.g. hinokitiol). These ingredients are claimed by the manufacturers to fight plaque, bad breath and prevent gum disease. A 2020 systematic metareview found that herbal toothpastes are as effective as non-herbal toothpastes in reducing dental plaque at shorter period of follow-up (4 weeks). However, this evidence comes from low-quality studies.
The stannous (tin) ion, commonly added to toothpastes as stannous fluoride or stannous chloride, has been shown to have antibacterial effects in the mouth. Research has shown that stannous fluoride-containing toothpaste inhibits extracellular polysaccharide (EPS) production in a multispecies biofilm greater than sodium fluoride-containing toothpaste. This is thought to contribute to a reduction in plaque and gingivitis when using stannous fluoride-containing toothpastes when compared to other toothpastes, and has been evidenced through numerous clinical trials. In addition to its antibacterial properties, stabilised stannous fluoride toothpastes have been shown to protect against dental erosion and dentine hypersensitivity, making it a multifunctional component in toothpaste formulations.
Flavorants
Toothpaste comes in a variety of colors and flavors, intended to encourage use of the product. The three most common flavorants are peppermint, spearmint, and wintergreen. Toothpaste flavored with peppermint-anise oil is popular in the Mediterranean region. These flavors are provided by the respective oils, e.g. peppermint oil. More exotic flavors include Anethole anise, apricot, bubblegum, cinnamon, fennel, lavender, neem, ginger, vanilla, lemon, orange, and pine. Alternatively, unflavored toothpastes exist.
Remineralizing agents
Chemical repair (remineralization) of early tooth decay is promoted naturally by saliva. However, this process can be enhanced by various remineralisation agents. Fluoride promotes remineralization, but is limited by bioavailable calcium. Casein phosphopeptide stabilised amorphous calcium phosphate (CPP-ACP) is a toothpaste ingredient containing bioavailable calcium that has been widely researched to be the most clinically effective remineralization agent that enhances the action of saliva and fluoride. Peptide-based systems, hydroxyapatite nanocrystals and a variety of calcium phosphates have been advocated as remineralization agents; however, more clinical evidence is required to substantiate their effectiveness.
Miscellaneous components
Agents are added to suppress the tendency of toothpaste to dry into a powder. Included are various sugar alcohols, such as glycerol, sorbitol, or xylitol, or related derivatives, such as 1,2-propylene glycol and polyethyleneglycol. Strontium chloride or potassium nitrate is included in some toothpastes to reduce sensitivity. Two systemic meta-analysis reviews reported that arginine, and calcium sodium phosphosilicate – CSPS containing toothpastes are also effective in alleviating dentinal hypersensitivity respectively. Another randomized clinical trial found superior effects when both formulas were combined.
Sodium polyphosphate is added to minimize the formation of tartar.
Chlorohexidine mouthwash has been popular for its positive effect on controlling plaque and gingivitis, however, a systemic review studied the effects of Chlorhexidine toothpastes and found insufficient evidence to support its use, tooth surface discoloration was observed as a side effect upon using it, which is considered a negative side effect that can affect patients' compliance.
Sodium hydroxide, also known as lye or caustic soda, is listed as an inactive ingredient in some toothpaste, for example Colgate Total.
Xylitol
A systematic review reported two out of ten studies by the same authors on the same population showed toothpastes with xylitol as an ingredient were more effective at preventing dental caries in permanent teeth of children than toothpastes containing fluoride alone. Furthermore, xylitol has not been found to cause any harmful effects. However, further investigation into the efficacy of toothpastes containing xylitol is required as the currently available studies are of low quality and high risk of bias.
Safety
Fluoride
Fluoride-containing toothpaste can be acutely toxic if swallowed in large amounts, but instances are exceedingly rare and result from prolonged and excessive use of toothpaste (i.e. several tubes per week). Approximately 15 mg/kg body weight is the acute lethal dose, even though as small amount as 5 mg/kg may be fatal to some children.
The risk of using fluoride is low enough that the use of full-strength toothpaste (1350–1500 ppm fluoride) is advised for all ages. However, smaller volumes are used for young children, for example, a smear of toothpaste until three years old. A major concern of dental fluorosis is for children under 12 months ingesting excessive fluoride through toothpaste. Nausea and vomiting are also problems which might arise with topical fluoride ingestion.
Diethylene glycol
The inclusion of sweet-tasting but toxic diethylene glycol in Chinese-made toothpaste led to a recall in 2007 involving multiple toothpaste brands in several nations. The world outcry made Chinese officials ban the practice of using diethylene glycol in toothpaste.
Triclosan
Reports have suggested triclosan, an active ingredient in many kinds of toothpastes, can combine with chlorine in tap water to form chloroform. An animal study revealed the chemical might modify hormone regulation, and many other lab researches proved bacteria might be able to develop resistance to triclosan in a way which can help them to resist antibiotics also.
Polyethylene glycol – PEG
PEG is a common ingredient in some of the formulas of toothpastes; it is a hydrophilic polymer that acts as a dispersant in toothpastes. Also, it is used in many cosmetic and pharmaceutical formulas, for example: ointments, osmotic laxatives, some of the nonsteroidal anti-inflammatory drugs, other medications and household products. However, 37 cases of PEG hypersensitivity (delayed and immediate) to PEG-containing substances have been reported since 1977, suggesting that they have unrecognized allergenic potential.
Miscellaneous issues and debates
With the exception of toothpaste intended to be used on pets such as dogs and cats, and toothpaste used by astronauts, most toothpaste is not intended to be swallowed, and doing so may cause nausea or diarrhea. Tartar fighting toothpastes have been debated. Sodium lauryl sulfate (SLS) has been proposed to increase the frequency of mouth ulcers in some people, as it can dry out the protective layer of oral tissues, causing the underlying tissues to become damaged. In studies conducted by the university of Oslo on recurrent aphthous ulcers, it was found that SLS has a denaturing effect on the oral mucin layer, with high affinity for proteins, thereby increasing epithelial permeability. In a double-blind cross-over study, a significantly higher frequency of aphthous ulcers was demonstrated when patients brushed with an SLS-containing versus a detergent-free toothpaste. Also patients with Oral Lichen Planus who avoided SLS-containing toothpaste benefited.
Alteration of taste perception
After using toothpaste, orange juice and other fruit juices are known to have an unpleasant taste if consumed shortly afterwards. Sodium lauryl sulfate, used as a surfactant in toothpaste, alters taste perception. It can break down phospholipids that inhibit taste receptors for sweetness, giving food a bitter taste. In contrast, apples are known to taste more pleasant after using toothpaste. Distinguishing between the hypotheses that the bitter taste of orange juice results from stannous fluoride or from sodium lauryl sulfate is still an unresolved issue and it is thought that the menthol added for flavor may also take part in the alteration of taste perception when binding to lingual cold receptors.
Whitening toothpastes
Many toothpastes make whitening claims. Abrasion is the principal way that toothpaste removes stains, and toothpastes that are not marketed as whitening can still remove stains by abrasion. Some of these toothpastes contain peroxide, the same ingredient found in tooth bleaching gels. Whitening toothpaste cannot alter the natural color of teeth or reverse discoloration by penetrating surface stains or decay. To remove surface stains, whitening toothpaste may include abrasives to gently polish the teeth or additives such as sodium tripolyphosphate to break down or dissolve stains. When used twice a day, whitening toothpaste typically takes two to four weeks to make teeth appear whiter. Whitening toothpaste is generally safe for daily use, but excessive use might damage tooth enamel. A recent systematic review in 2017 concluded that nearly all dentifrices that are specifically formulated for tooth whitening were shown to have a beneficial effect in reducing extrinsic stains, irrespective of whether or not a chemical discoloration agent was added. However, the whitening process can permanently reduce the strength of the teeth, as the process scrapes away a protective outer layer of enamel.
Herbal and natural toothpastes
Herbal toothpastes are marketed to consumers who wish to avoid some of the artificial ingredients commonly found in regular toothpastes. The ingredients found in so-called natural toothpastes vary widely but often include baking soda, aloe, eucalyptus oil, myrrh, camomile, calendula, neem, toothbrush tree, plant extract (strawberry extract), and essential oils. Many herbal toothpastes do not contain fluoride or sodium lauryl sulfate.
A 2020 meta-analysis showed some evidence for the efficacy of herbal toothpaste, albeit from poor quality studies. According to a study by the Delhi Institute of Pharmaceutical Sciences and Research, many of the herbal toothpastes being sold in India were adulterated with nicotine.
Charcoal has also been incorporated in toothpaste formulas; however, there is no evidence to determine its safety and effectiveness, and the American Dental Association does not recommend its use.
Government regulation
In the United States toothpaste is regulated by the U.S. Food and Drug Administration as a cosmetic, except for ingredients with a medical purpose, such as fluoride, which are regulated as drugs. Drugs require scientific studies and FDA approval in order to be legally marketed in the United States, but cosmetic ingredients do not require pre-approval, except for color additives. The FDA does have labelling and requirements and bans certain ingredients.
Striped toothpaste
Striped toothpaste was invented by Leonard Marraffino in 1955. The patent (US patent , issued 1957) was subsequently sold to Unilever, which marketed the novelty under the Stripe brand-name in the early 1960s. This was followed by the introduction of the Signal brand in Europe in 1965 (UK patent 813,514). Although Stripe was initially very successful, it never again achieved the 8% market share that it cornered during its second year.
Marraffino's design, which remains in use for single-color stripes, is simple. The main material, usually white, sits at the crimp end of the toothpaste tube and makes up most of its bulk. A thin pipe, through which that carrier material will flow, descends from the nozzle to it. The stripe-material (this was red in Stripe) fills the gap between the carrier material and the top of the tube. The two materials are not in separate compartments, but they are sufficiently viscous that they will not mix. When pressure is applied to the toothpaste tube, the main material squeezes down the thin pipe to the nozzle. Simultaneously, the pressure applied to the main material causes pressure to be forwarded to the stripe material, which thereby issues out through small holes (in the side of the pipe) onto the main carrier material as it is passing those holes.
In 1990, Colgate-Palmolive was granted a patent (USPTO ) for two differently colored stripes. In this scheme, the inner pipe has a cone-shaped plastic guard around it, and about halfway up its length. Between the guard and the nozzle-end of the tube is a space for the material for one color, which issues out of holes in the pipe. On the other side of the guard is space for second stripe-material, which has its own set of holes.
In 2016, Colgate-Palmolive was granted a patent (USPTO ) for suitable sorts of differently colored toothpastes to be filled directly into tubes to produce a striped mix without any separate compartments. This required adjustment of the different components' behavior (rheology) so that stripes are produced when the tube is squeezed.
Striped toothpaste should not be confused with layered toothpaste. Layered toothpaste requires a multi-chamber design (e.g. USPTO ), in which two or three layers extrude out of the nozzle. This scheme, like that of pump dispensers (USPTO ), is more complicated (and thus, more expensive to manufacture) than either the Marraffino design or the Colgate designs.
The iconic depiction of a wave-shaped blob of toothpaste sitting on a toothbrush is called a "nurdle".
Tooth powder
Tooth powders for use with toothbrushes came into general use in the 19th century in Britain. Most were homemade, with chalk, pulverized brick, or salt as ingredients. An 1866 Home Encyclopedia recommended pulverized charcoal, and cautioned that many patented tooth powders that were commercially marketed did more harm than good.
Arm & Hammer marketed a baking soda-based toothpowder in the United States until approximately 2000, and Colgate currently markets toothpowder in India and other countries.
Modern toothpaste
An 18th-century American and British toothpaste recipe called for burned bread. Another formula around this time called for dragon's blood (a resin), cinnamon, and burned alum.
In 1873 the Colgate company began the mass production of aromatic toothpaste in jars.
By 1900, a paste made of hydrogen peroxide and baking soda was recommended for use with toothbrushes. Pre-mixed toothpastes were first marketed in the 19th century, but did not surpass the popularity of tooth-powder until World War I.
Together with Willoughby D. Miller, Newell Sill Jenkins developed the first toothpaste containing disinfectants, branded as Kolynos. The name is a combination of two Greek words, meaning "beautifier" and "disease preventer". Numerous attempts to produce the toothpaste by pharmacists in Europe proved uneconomic. After returning to the US, he continued experimenting with Harry Ward Foote (1875–1942), professor of chemistry at Sheffield Chemical Laboratory of Yale University. After 17 years of development of Kolynos and clinical trials, Jenkins retired and transferred the production and distribution to his son Leonard A. Jenkins, who brought the first toothpaste tubes on the market on April 13, 1908. Within a few years the company expanded in North America, Latin America, Europe and the Far East. A branch operation opened in London in 1909. In 1937, Kolynos was produced in 22 countries and sold in 88 countries. Kolynos has been sold mainly in South America and in Hungary. Colgate-Palmolive took over the production of American Home Products in 1995 at a cost of one billion US dollars.
Fluoride was first added to toothpastes in the 1890s. Tanagra, containing calcium fluoride as the active ingredient, was sold by Karl F. Toellner Company, of Bremen, Germany, based upon the early work of chemist Albert Deninger. An analogous invention by Roy Cross, of Kansas City, Missouri, was initially criticized by the American Dental Association (ADA) in 1937. Fluoride toothpastes developed in the 1950s received the ADA's approval. To develop the first ADA-approved fluoride toothpaste, Procter & Gamble started a research program in the early 1940s. In 1950, Procter & Gamble developed a joint research project team headed by Joseph C. Muhler at Indiana University to study new toothpaste with fluoride. In 1955, Procter & Gamble's Crest launched its first clinically proven fluoride-containing toothpaste. On August 1, 1960, the ADA reported that "Crest has been shown to be an effective anticavity (decay preventative) dentifrice that can be of significant value when used in a conscientiously applied program of oral hygiene and regular professional care."
In 1980, the Japanese company, Sangi Co., Ltd., launched APADENT, the world's first remineralizing toothpaste to use a nano-form of hydroxyapatite, the main component of tooth enamel, rather than fluoride, to remineralize areas of mineral loss below the surface of tooth enamel (incipient caries lesions). After many years of laboratory experiments and field trials, its hydroxyapatite ingredient was approved as an active anti-caries agent by the Japanese Ministry of Health in 1993, and given the name Medical Hydroxyapatite to distinguish it from other forms of hydroxyapatite used in toothpaste, such as dental abrasives.
In 2006, BioRepair appeared in Europe with the first European toothpaste containing synthetic hydroxylapatite as an alternative to fluoride for the remineralization and reparation of tooth enamel. The "biomimetic hydroxylapatite" is intended to protect the teeth by creating a new layer of synthetic enamel around the tooth instead of hardening the existing layer with fluoride that chemically changes it into fluorapatite.
Dispensing
Toothpaste is usually dispensed via a collapsible tube or with a more rigid pump. Several traditional and innovative designs have been developed. The dispenser must be matched to the flow properties of the toothpaste.
In 1880, Doctor Washington Sheffield of New London, CT manufactured toothpaste into a collapsible tube, Dr. Sheffield's Creme Dentifrice. He had the idea after his son traveled to Paris and saw painters using paint from tubes. In York in 1896, Colgate-Palmolive Dental Cream was packaged in collapsible tubes imitating Sheffield. The original collapsible toothpaste tubes were made of lead.
| Biology and health sciences | Hygiene products | Health |
870898 | https://en.wikipedia.org/wiki/Firing%20pin | Firing pin | A firing pin or striker is a part of the firing mechanism of a firearm that impacts the primer in the base of a cartridge and causes it to fire. In firearms terminology, a striker is a particular type of firing pin where a compressed spring acts directly on the firing pin to provide the impact force rather than it being struck by a hammer.
The terms may also be used for a component of equipment or a device which has a similar function. Such equipment or devices include: artillery, munitions and pyrotechnics.
Firearms
The typical firing pin is a thin, simple rod with a hardened, rounded tip that strikes and crushes the primer. The rounded end ensures the primer is indented rather than pierced (to contain propellant gasses). It sits within a hole through the breechblock and is struck by the hammer when the trigger is "pulled". A light firing-pin spring is often used to keep the firing pin rearward. It may be termed a firing-pin return spring, since it returns it to the unfired position. In semi-automatic firearms, this prevents premature firing from the inertia of the firing pin as the breech mechanism closes in the reloading part of the firing cycle. Firing pins of this type are often too short to contact the primer when the hammer is resting against it. This is a safety measure to prevent discharge from external forces such as a drop. Firing relies upon the transfer of momentum from the hammer to give the firing pin sufficient impact energy to cause firing.
Type of cartridge
The two main types of metallic cartridges used in modern firearms are centerfire and rimfire. In centerfire cartridges, the primer is located in the center of the base of the cartridge. Firing pins for centerfire cartridges usually have a round cross-section and their movement is usually through a hole in the breechblock along the axis of the center of the barrel's bore.
Rimfire cartridges however, must be struck on the base about the rim of the cartridge. While rimfire firearms may use a firing pin with a round cross-section, it is common for them to be flat, with a square or rectangular cross-section and a blunt chisel point. Flat firing pins can be stamped from flat metal stock and usually operate in a slot cut in the breechblock (rather than a hole) that is parallel but offset from the centerline of the barrel. These production methods are generally simpler and reduce production costs. It is generally recommended not to excessively "dry fire" rimfire firearms (ie firing without a chambered round) as it is possible for the firing pin to strike the face of the camber and deform it or damage the firing pin.
Historical cartridges
In 1808 by the Swiss gunsmith Jean Samuel Pauly in association with French gunsmith François Prélat created the first cartridges to integrate a primer and be self-contained. The paper cartridge used a metal base with a through-hole coated in a percussive priming compound. Pauly also developed a breech-loading shotgun for his cartridge, using a firing pin and external hammer.
The Dreyse needle gun of 1836 uses a paper cartridge with a priming as part of a sabot which cradles the projectile and is forward of the propelling charge. The needle-like firing pin projects from the bolt-face and pierces the cartridge when the breech is closed. On firing, the spring-loaded needle strikes the priming in the sabot. Unlike Pauly's cartridge, which was not widely accepted, Dreyse's rifle was adopted by Prussia as its infantry service rifle. It was the first military breechloader to use a self-contained cartridge and consequently, the first to employ a firing pin.
The pinfire cartridge patented in 1835, uses a metallic cartridge with an integrated firing pin located radially near the base of the cartridge. The pin needs to be aligned with a corresponding slot in the chamber; a disadvantage compared with rimfire and centerfire cartridges that followed and that are also safer.
Side-lock hammers
Early rifle designs that fired metallic cartridges typically used a side-lock mechanism, with the hammer mounted to one side rather than inline with the axis of the barrel. In the trapdoor Springfield Model 1865 (and similar) the rear of the firing pin tube within the breechblock is angled away from the centerline of the barrel toward the hammer. The Sharps rifle uses a firing pin block to solve this alignment problem. The block sits within a recess in the breechblock. When struck by the hammer, the whole block is propelled forward. That part of the block with the firing pin sits on the centerline of the barrel and strikes the primer.
Fixed firing pins
Many revolvers use a firing pin that is fixed to the hammer. Simple blowback sub-machine guns that fire from the open-bolt position often have a fixed firing pin that protrudes from the face of the bolt. As the bolt fully closes on the breech the primer of the newly chambered round is struck, causing the cartridge to fire. The Owen and F1 submachine gun are examples that use bolt-face fixed firing-pins. Some mortars use a fixed firing pin mounted in the breech plug. When a mortar round is dropped down the barrel, a primer in the base of the mortar round strikes the firing pin and ignites the propelling charge.
Floating firing pins
In firearms terminology, a floating firing pin is one which is unrestricted by a firing-pin return spring or similar. While it will be captive and unable to simply fall out, either forward or backward, it is otherwise free to slide within these stops. The trapdoor Springfield Model 1865 is an example of a floating firing pin.
Striker
Hammer-operated firing mechanisms use a relatively light firing pin and rely on a transfer of momentum received from a spring-loaded hammer, in the same fashion as a punch or chisel relays the blow from a mallet. In firearms terminology, a striker (or striker mechanism) derives the impact force to strike the primer from a spring acting directly upon the firing pin – similar to a crossbow, where the striker (firing pin) is like the crossbow bolt (arrow). A striker mechanism is very common in bolt-action firearms but not to the exclusion of hammer-operated mechanisms. It is also found in other actions where the breechblock reciprocates directly inline with the axis of the barrel.
A striker mechanism will consist of the striker spring (firing spring) and the striker. The striker spring is a relatively strong spring sufficient to initiate firing. A typical striker consists of a narrow striking point, a heavier section that acts as a spring guide for the striker spring, a shoulder to restrain the spring, and a catch piece which is engaged by the trigger sear to hold the spring under tension when "cocked" and ready to fire. The striker spring is compressed between the striker's shoulder and the rear of the breechblock. A striker may be assembled from several component; however, the stored energy in the striker spring is transferred directly to the striker and then to the primer without any intermediate transfer of energy or momentum. As striker mechanisms combine both functions of hammer and firing pin in one piece, they are generally considered to be mechanically simpler but are more robust in construction than a typical firing pin.
Many small-caliber rimfire bolt-action rifles and some centerfire automatic weapons (e. g., vz. 58) may appear to have a striker-operated firing mechanism but are actually a type of linear hammer. The hammer can be likened to that of a pile driver and is mainly contained within the bolt. It is much like the striker already described except that the "hammer" upon which the firing spring acts and the firing pin are separate units. Confusingly, parts lists will often refer to this type of hammer as a "striker".
Striker-fired (or similar) bolt action firearms may be classified as cock-on-close or cock-on-open.
Cock-on-close
When the breech is opened and retracted rearward, the striker is also carried rearward so that the striker catch passes over the trigger sear. When the bolt is pushed forward to close the breech, the striker catch is held by the trigger sear. The firer must close the bolt with sufficient force to overcome the force exerted by the cocking spring. Notably, the Lee–Enfield and Belgian Mauser cock on closing as do many small-caliber rimfire bolt-action rifles.
Cock-on-open
The breech of a bolt action rifle is opened by first rotating the bolt handle. In cock-on-open operation, this rotation acts on a cam (similar to the action of a screw thread) which retracts the striker, compressing the cocking spring and holding it there. When the cocking handle is rotated closed, the cocking cam disengages but the striker is retained in the cocked position by the trigger sear.
Introduced in the Mauser Model 1871, it significantly reduced the risk of accidental discharge upon closing. The system of operation was widely adopted and is used almost exclusively in modern centerfire rifle designs. The Mauser Gewehr 98, the Mosin–Nagant and M1903 Springfield are examples of service rifles using this type of operation.
Other applications
Modern era guns in general (and not just firearms) use a firing pin of some description to initiate firing. Mechanical contact fuzes in explosive ordnance will employ a firing pin or striker to initiate detonation. Such devices include: artillery projectiles, aerial bombs and land mines. In landmines, non-metallic firing pins, made from ceramics for example, may be used to minimise their magnetic signature.
The M127A1 signal rocket (and other similar flares) have a primer in the base of the disposable launch tube. The cap contains a fixed firing-pin inside. The flare is fired by placing the cap over the base and striking it by hand.
Hand grenades of the type that use a safety lever (such as the M26 grenade) use a striker that is similar to the classic spring-loaded mousetrap. It is held under tension until the lever is released and then flips over to strike the primer cap. Some chemical oxygen generators use a primer and mousetrap type striker to initiate the chemical reaction.
Images
| Technology | Mechanisms_2 | null |
871210 | https://en.wikipedia.org/wiki/Utricularia | Utricularia | Utricularia, commonly and collectively called the bladderworts, is a genus of carnivorous plants consisting of approximately 233 species (precise counts differ based on classification opinions; a 2001 publication lists 215 species). They occur in fresh water and wet soil as terrestrial or aquatic species across every continent except Antarctica. Utricularia are cultivated for their flowers, which are often compared with those of snapdragons and orchids, especially amongst carnivorous plant enthusiasts.
All Utricularia are carnivorous and capture small organisms by means of bladder-like traps. Terrestrial species tend to have tiny traps that feed on minute prey such as protozoa and rotifers swimming in water-saturated soil. The traps can range in size from . Aquatic species, such as U. vulgaris (common bladderwort), possess bladders that are usually larger and can feed on more substantial prey such as water fleas (Daphnia), nematodes and even fish fry, mosquito larvae and young tadpoles. Despite their small size, the traps are extremely sophisticated. In the active traps of the aquatic species, prey brush against trigger hairs connected to the trapdoor. The bladder, when "set", is under negative pressure in relation to its environment so that when the trapdoor is mechanically triggered, the prey, along with the water surrounding it, is sucked into the bladder. Once the bladder is full of water, the door closes again, the whole process taking only ten to fifteen milliseconds.
Bladderworts are unusual and highly specialized plants, and the vegetative organs are not clearly separated into roots, leaves, and stems as in most other angiosperms. Utricularia lack a root system.
Bladder traps are recognized as one of the most sophisticated structures in the plant kingdom.
Description
The main part of a bladderwort plant always lies beneath the surface of its substrate. Terrestrial species sometimes produce a few photosynthetic leaf-shoots. The aquatic species can be observed below the surfaces of ponds and streams.
Plant structure
Most species form long, thin, sometimes branching stems or stolons beneath the surface of their substrate, whether that be pond water or dripping moss in the canopy of a tropical rainforest. To these stolons are attached both the bladder traps and photosynthetic leaf-shoots, and in terrestrial species the shoots are thrust upward through the soil into the air or along the surface.
The name bladderwort refers to the bladder-like traps. The aquatic members of the genus have the largest and most obvious bladders, and these were initially thought to be flotation devices before their carnivorous nature was discovered.
Etymology
The generic name Utricularia is derived from the Latin utriculus, a word which has many related meanings but which most commonly means wine flask, leather bottle or bagpipe.
Flowers
Flowers are the only part of the plant clear of the underlying soil or water. They are usually produced at the end of thin, often vertical inflorescences. They can range in size from wide, and have two asymmetric labiate (unequal, lip-like) petals, the lower usually significantly larger than the upper. They can be of any colour, or of many colours, and are similar in structure to the flowers of a related carnivorous genus, Pinguicula.
The flowers of aquatic varieties like U. vulgaris are often described as similar to small yellow snapdragons, and the Australian species U. dichotoma can produce the effect of a field full of violets on nodding stems. The epiphytic species of South America, however, are generally considered to have the showiest, as well as the largest, flowers. It is these species that are frequently compared with orchids.
Certain plants in particular seasons might produce closed, self-pollinating (cleistogamous) flowers; but the same plant or species might produce open, insect-pollinated flowers elsewhere or at a different time of year, and with no obvious pattern. Sometimes, individual plants have both types of flower at the same time: aquatic species such as U. dimorphantha and U. geminiscapa, for example, usually have open flowers riding clear of the water and one or more closed, self-pollinating flowers beneath the water. Seeds are numerous and small and for the majority of species are long.
Distribution and habitat
Utricularia can survive almost anywhere where there is fresh water for at least part of the year; only Antarctica and some oceanic islands have no native species. The greatest species diversity for the genus is seen in South America, with Australia coming a close second. In common with most carnivorous plants, they grow in moist soils which are poor in dissolved minerals, where their carnivorous nature gives them a competitive advantage; terrestrial varieties of Utricularia can frequently be found alongside representatives of the carnivorous genera–Sarracenia, Drosera and others–in very wet areas where continuously moving water removes most soluble minerals from the soil.
Utricularia has a variety of life forms, including terrestrial, lithophytic, aquatic, epiphytic, and rheophytic forms which are all highly adapted for their environments.
About 80% of the species are terrestrial, and most inhabit waterlogged or wet soils, where their tiny bladders can be permanently exposed to water in the substrate. Frequently they will be found in marshy areas where the water table is very close to the surface. Most of the terrestrial species are tropical, although they occur worldwide.
Approximately 20% of the species are aquatic. Most of these drift freely over the surface of ponds and other still, muddy-bottomed waters and only protrude above the surface when flowering, although a few species are lithophytic and adapted to rapidly moving streams or even waterfalls. The plants are usually found in acidic waters, but they are quite capable of growing in alkaline waters and would very likely do so were it not for the higher level of competition from other plants in such areas. Aquatic Utricularia are often split into two categories: suspended and affixed aquatic. Suspended aquatics are species which are not rooted into the ground and are free-floating, often found in nutrient poor sites. Conversely, fixed aquatics are species which have at least some of their shoots rooted into the ground. These plants often have dimorphic shoots, some which are leafy, green, and often bladderless which float in the water, and others which are white and coated with bladders that affix the plant to the ground. Utricularia vulgaris is an aquatic species and grows into branching rafts with individual stolons up to one metre or longer in ponds and ditches throughout Eurasia.
Some South American tropical species are epiphytes, and can be found growing in wet moss and spongy bark on trees in rainforests, or even in the watery leaf-rosettes of other epiphytes such as various Tillandsia (a type of bromeliad) species. Epiphytic Utricularia are often known for their orchid-like flowers and are the most ornamentally sought after. Rosette-forming epiphytes such as U. nelumbifolia put out runners, searching for other nearby bromeliads to colonise.
There are also a few lithophytic species which live on wet surfaces of cliffs and mossy rocks and rheophytic species which live in shallow rivers and streams.
The plants are as highly adapted in their methods of surviving seasonally inclement conditions as they are in their structure and feeding habits. Temperate perennials can require a winter period in which they die back each year, and they will weaken in cultivation if they are not given it; tropical and warm-temperate species, on the other hand, require no dormancy. Floating bladderworts in cold temperate zones such as the UK and Siberia can produce winter buds called turions at the extremities of their stems: as the autumnal light fails and growth slows down, the main plant may rot away or be killed by freezing conditions, but the turions will separate and sink to the bottom of the pond to rest beneath the coming ice until the spring, when they will return to the surface and resume growth. Many Australian species will grow only during the wet season, reducing themselves to tubers only long to wait out the dry season. Other species are annual, returning from seed each year.
Dispersal and life form evolution
The ancestral line of Utricularia is thought to have been terrestrial. From terrestrial forms, epiphytic forms evolved independently three times and aquatic life forms arose four times in genus Utricularia.
Biogeographic patterns associated with the boreotropic hypothesis lists the origin of Lentibulariaceae to temperate Eurasia or tropical America. Based on fossilised pollen and insular separation, the last common ancestor of Genlisea-Utricularia clade was found to be a South American lineage that arose 39 mya. Utricularia probably diverged from its sister genus 30 mya and subsequently dispersed to Australia, represented by subgenus Polypompholyx, and to Africa. There were most likely other transcontinental dispersals, one of which is represented by sect. Nelipus. The colonization of Utricularia to North America probably occurred 12mya from South America. The dispersal of Utricularia to Eurasia probably occurred through the Bering Strait via long-distance dispersal 4.7 mya.
Carnivory
Physical description of the trap
Authorities on the genus, such as botanists Peter Taylor and Francis Ernest Lloyd, agree that the vacuum-driven bladders of Utricularia are the most sophisticated carnivorous trapping mechanism to be found anywhere in the plant kingdom. The bladders are usually shaped similarly to broad beans (though they come in various shapes) and attach to the submerged stolons by slender stalks.
Bladders are hollow underwater suction cups, also known as utricles, that possess a valve with bristles that open and close.
The bladder walls are very thin and transparent but are sufficiently inflexible to maintain the bladder's shape despite the vacuum created within. The entrance, or 'mouth', of the trap is a circular or oval flap whose upper half is joined to the body of the trap by very flexible, yielding cells which form an effective hinge. The door rests on a platform formed by the thickening of the bladder wall immediately underneath. A soft but substantial membrane called the velum stretches in a curve around the middle of this platform, and helps seal the door. A second band of springy cells crosses the door just above its lower edge and provides the flexibility for the bottom of the door to become a bendable 'lip' which can make a perfect seal with the velum.
The outer cells of the whole trap excrete mucilage and under the door, this is produced in greater quantities and contains sugars. The mucilage certainly contributes towards the seal, and the sugars may help to attract prey.
Terrestrial species, like U. sandersonii have tiny traps (sometimes as small as 0.2 mm; 1/100") with a broad beak-like structure extending and curving down over the entrance; this forms a passageway to the trapdoor and may help prevent the trapping and ingestion of inorganic particles. Aquatic species, like U. inflata tend to have larger bladders—up to —and the mouth of the trap is usually surrounded not by a beak but by branching antennae, which serve both to guide prey animals to the trap's entrance and to fend the trap mouth away from larger bodies which might trigger the mechanism needlessly. Epiphytic species have unbranched antennae which curve in front of the mouth and probably serve the same purpose, although it has been observed that they are also capable of holding a pocket of water in front of the mouth by capillary action, and that this assists with the trapping action.
Trapping mechanism
The trapping mechanism of Utricularia is purely mechanical; no reaction from the plant (irritability) is required in the presence of prey, in contrast with the triggered mechanisms employed by Venus flytraps (Dionaea), waterwheels (Aldrovanda), and many sundews (Drosera). The only active mechanism involved is the constant pumping out of water through the bladder walls by active transport.
As water is pumped out, the bladder's walls are sucked inwards by the negative pressure created, and any dissolved material inside the bladder becomes more concentrated. The sides of the bladder bend inwards, storing potential energy like a spring. Eventually, no more water can be extracted, and the bladder trap is 'fully set' (technically, osmotic pressure rather than physical pressure is the limiting factor).
Extending outwards from the bottom of the trapdoor are several long bristle-stiff protuberances that are sometimes referred to as trigger hairs or antennae but which have no similarity to the sensitive triggers found in Dionaea and Aldrovanda. In fact, these bristles are simply levers. The suction force exerted by the primed bladder on the door is resisted by the adhesion of its flexible bottom against the soft-sealing velum. The equilibrium depends quite literally on a hair trigger, and the slightest touch to one of the lever hairs will deform the flexible door lip enough to create a tiny gap, breaking the seal.
Once the seal is disturbed, the bladder walls instantly spring back to a more rounded shape; the door flies open and a column of water is sucked into the bladder. The animal which touched the lever, if small enough, is inevitably drawn in, and as soon as the trap is filled, the door resumes its closed position—the whole operation being completed in as little as one-hundredth of a second.
Once inside, the prey is dissolved by digestive secretions. This generally occurs within a few hours, although some protozoa appear to be highly resistant and have been observed to live for several days inside the trap. All the time, the trap walls continue to pump out water, and the bladder can be ready for its next capture in as little as 15 to 30 minutes.
Microbial relationships
The bladders of Utricularia often culture a mutualistic community of microbes, which may be a very important factor in digestion of prey within Utricularia. Bacteria consume dissolved organic material which is not able to be directly ingested by larger organisms. When bacteria absorb dissolved organic material, they also release nutrients, which facilitates photo-autotrophic growth. As Utricularias trap is sealed and contains all the needed components of a microbial food web, one can assume that much enzyme activity and available nutrients in Utricularias trap fluid are derived from these microbial communities. Additionally, Utricularia traps often collect a diversity of microplankton and detritus. When this periphyton is dissolved into basic nutrients within the bladder environment, bacterial enzymes help aid in digestion. Therefore, carbon secretion and periphyton utilization in the utricles enable Utricularia to live with relatively little competition.
Mutualism could have been an important association in aquatic Utricularia trap evolution as these microbes may have allowed these plants to acquire the needed nutrients when they lost their roots, as they may have had issues acquiring phosphorus. Phosphorus was found to be the most important factor in Utricularia nutrition, which helps explain why Utricularia bladders are found with a wide diversity of bacteria to aid in phosphorus digestion.
Enhanced respiration
Utricularia have significantly greater respiration rates than most vegetative tissue, primarily due to their complex energy-dependent traps. Upon triggering, prey is captured through a two-step ATP-driven ion-pumping process where organisms are sucked in by internal negative pressure achieved by pumping water out of the trap and into the external environment. Recent research suggests that COX subunit I (COX1), a rate limiting enzyme in the cellular respiration pathway associated with the synthesis of ATP, has evolved under positive Darwinian selection in the Utricularia–Genlisea clade. There appear to be adaptive substitutions of two contiguous cysteines (C-C motif) at the docking point of COX1 helix 3 and cytochrome c. This C-C motif, absent in ~99.9% of databased Eukaryota, Archaea, and Bacteria, suggests a conformational change that might decouple electron transport from proton pumping. By doing so, the intermembrane space could sequester protons are store them until the ATP is needed. Such decoupling would allow Utricularia to optimize power output (energy × rate) during times of need, albeit with a 20% cost in energy efficiency.
According to the ROS mutation hypothesis, the sequestration of these protons has cellular consequences, which could lead to nucleotide substitutions. Oxidative phosphorylation is an imperfect process, which allows electrons to leak into the lumen, and only partially reduce oxygen. This partially reduced oxygen is a reactive oxygen species (ROS) which can be very harmful, unlike its fully reduced counterpart, the water molecule. When there is greater potential change between the lumen and intermembrane space, the leakiness of the electron transport chain also increases, therefore creating a higher production of ROS in the mitochondria of Utricularia. ROS is harmful to cells, as it produces damage to nucleotides and helical DNA. Therefore, the increased cellular respiration of Utricularia bladders combined with the unique sequestration of protons could lead to its high nucleotide substitution rates, and therefore its wide diversity.
This structural evolution seems highly unlikely to have arisen by chance alone; therefore, many researchers suggest this key adaption in Utricularia allowed for radical morphological evolution of relatively simple trap structures to highly complex and efficient snares. This adaptation may have enhanced the genus' fitness by increasing its range of prey, rate of capture, and retention of nutrients during prey decomposition.
Lloyd's experiments
In the 1940s, Francis Ernest Lloyd conducted extensive experiments with carnivorous plants, including Utricularia, and settled many points which had previously been the subject of conjecture. He proved that the mechanism of the trap was purely mechanical by both killing the trigger hairs with iodine and subsequently showing that the response was unaffected, and by demonstrating that the trap could be made ready to spring a second (or third) time immediately after being set off if the bladder's excretion of water were helped by a gentle squeeze; in other words, the delay of at least fifteen minutes between trap springings is due solely to the time needed to excrete water, and the triggers need no time to recover irritability (unlike the reactive trigger hairs of Venus Flytraps, for example).
He tested the role of the velum by showing that the trap will never set if small cuts are made to it; and showed that the excretion of water can be continued under all conditions likely to be found in the natural environment, but can be prevented by driving the osmotic pressure in the trap beyond normal limits by the introduction of glycerine.
Ingestion of larger prey
Lloyd devoted several studies to the possibility, often recounted but never previously accounted for under scientific conditions, that Utricularia can consume larger prey such as young tadpoles and mosquito larvae by catching them by the tail, and ingesting them bit by bit.
Prior to Lloyd, several authors had reported this phenomenon and had attempted to explain it by positing that creatures caught by the tail repeatedly set off the trap as they thrash about in an attempt to escape—even as their tails are actively digested by the plant. Lloyd, however, demonstrated that the plant is quite capable of ingestion by stages without the need of multiple stimuli.
He produced suitable artificial "prey" for his experiments by stirring albumen (egg white) into hot water and selecting shreds of an appropriate length and thickness. When caught by one end, the strand would gradually be drawn in, sometimes in sudden jumps, and at other times by a slow and continuous motion. Strands of albumen would often be fully ingested in as little as twenty minutes.
Mosquito larvae, caught by the tail, would be engulfed bit by bit. A typical example given by Lloyd showed that a larva of a size at the upper limit of what the trap could manage would be ingested stage by stage over the course of about twenty-four hours; but that the head, being rigid, would often prove too large for the mouth of the trap and would remain outside, plugging the door. When this happened, the trap evidently formed an effective seal with the head of the larva as it could still excrete water and become flattened, but it would nevertheless die within about ten days "evidently due to overfeeding".
Softer-bodied prey of the same size such as small tadpoles could be ingested completely, because they have no rigid parts and the head, although capable of plugging the door for a time, will soften and yield and finally be drawn in.
Very thin strands of albumen could be soft and fine enough to allow the trapdoor to close completely; these would not be drawn in any further unless the trigger hairs were indeed stimulated again. On the other hand, a human hair, finer still but relatively hard and unyielding, could prevent a seal being formed; these would prevent the trap from resetting at all due to leakage of water.
Lloyd concluded that the sucking action produced by the excretion of water from the bladder was sufficient to draw larger soft-bodied prey into the trap without the need for a second or further touch to the trigger levers. An animal long enough not to be fully engulfed upon first springing the trap, but thin and soft enough to allow the door to return fully to its set position, would indeed be left partly outside the trap until it or another body triggered the mechanism once again. However, the capture of hard bodies not fully drawn into the trap would prevent its further operation.
Genetics
Chris Whitewoods has developed a computational model of possible genetic regulation in Utricularia gibba to show how genes may control the formation of the upper and lower surfaces of flat leaves and how cup-shaped traps may have evolved from flat leaves. Changes in the gene expression of Utricularia can explain these structural changes. U. gibba leaves appear similar early in development but may develop into either a spherical trap or a cylindrical leaflet at later stages. Directional expansion of the leaf is suggested to be a crucial driver of the trap's morphogenesis. The upper and lower faces of the leaf are differentially associated with genetic markers. The marker UgPHV1 is associated with the upper leaf face. Trap primordia become spherical in shape, due to growth in both the longitudinal and transverse directions, when UgPHV1 / PHAVOLUTA (PHV) is restricted. Expression of UgPHV1 inhibits trap development and leads to the formation of leaflets. The same model can be used to describe shape development of other leaf shapes, including the pitcher-shaped Sarracenia trap, in terms of the spatial regulation of gene expression.
Increased respiration rates caused by mutated COXI may have caused two additional traits in the Utricularia–Genlisea clade: i) greatly increased rates of nucleotide substitution and ii) a dynamic decrease of genome size, including Utricularia species with some of the smallest haploid angiosperm genomes known. A recent study conducted three cDNA libraries from different organs of U. gibba (~80Mb) as part of a large scale Utricularia nuclear genome sequencing project. They recorded increased nucleotide substitution rates in chloroplast, mitochondrial, and cellular genomes. They also recorded increased levels of DNA repair-associated proteins and reactive oxygen species (ROS)-detox. ROS is a product of cellular metabolism that can potentially cause cellular damage when accumulated in high amounts. They determined the expression of DNA repair and ROS detox was ubiquitous rather than trap-specific. Due to this ubiquitous expression, relative ROS detoxification is expected to be lower in trap structures due to the high respiratory rate caused by trap activations, eventually leading to higher toxic effects and mutagenesis. Mutagenic action of enhanced ROS production may explain both high rates of nucleotide substitution and the dynamic evolution of genome size (via double strand breaks).
The dramatic shift in genome size and high mutation rates may have allowed for the variations observed in Utricularia bladder size, root structure, and relaxed body formation. Overall, the introduction of mutated COXI and high mutation rates provide a strong evolutionary hypothesis for the variability found in Utricularia species.
Species
Utricularia is the largest genus of carnivorous plants. It is one of the three genera that make up the Bladderwort family (Lentibulariaceae), along with the butterworts (Pinguicula) and corkscrew plants (Genlisea).
This genus was considered to have 250 species until Peter Taylor reduced the number to 214 in his exhaustive study The genus Utricularia – a taxonomic monograph, published by Her Majesty's Stationery Office in 1989. Taylor's classification is now generally accepted with modifications based on phylogenetic studies (see below).
The genus Polypompholyx, the pink petticoats, contained just two species of carnivorous plant, Polypompholyx tenella and Polypompholyx multifida, previously distinguished from the otherwise similar genus Utricularia by their possession of four calyx lobes rather than two. The genus has now been subsumed into Utricularia.
The genus Biovularia contained the species Biovularia olivacea (also known as B. brasiliensis or B. minima) and Biovularia cymbantha. The genus has been subsumed into Utricularia.
Phylogenetics
The following cladogram shows the relationship between various subgenera and sections. It summarizes the results of two studies (Jobson et al. 2003; Müller et al. 2004), following Müller et al. 2006. Since the sections Aranella and Vesiculina are polyphyletic, they show up multiple times in the cladogram (*). Some monotypic sections have not been included in the study, so that their place in this system is unclear. Sections that are not included below are Candollea, Chelidon, Choristothecae, Kamienskia, Martinia, Meionula, Mirabiles, Oliveria, Setiscapella, Sprucea, Steyermarkia, and Stylotheca in subgenus Utricularia; Minutae in subgenus Bivalvaria; and Tridentaria in subgenus Polypompholyx.
| Biology and health sciences | Lamiales | Plants |
871330 | https://en.wikipedia.org/wiki/Aldrovanda%20vesiculosa | Aldrovanda vesiculosa | Aldrovanda vesiculosa, commonly known as the waterwheel plant, is the sole extant species in the flowering plant genus Aldrovanda of the family Droseraceae. The plant captures small aquatic invertebrates using traps similar to those of the Venus flytrap. The traps are arranged in whorls around a central, free-floating stem, giving rise to the common name. This is one of the few plant species capable of rapid movement.
While the genus Aldrovanda is now monotypic, up to 19 extinct species are known in the fossil record. While the species displays a degree of morphological plasticity between populations, A. vesiculosa possesses a very low genetic diversity across its entire range.
A. vesiculosa has declined over the last century to only 50 confirmed extant populations worldwide. These are spread across Europe, Africa, Asia, and Australia. However, potentially invasive populations exist in the eastern United States. It is kept by hobbyists.
Morphology
Aldrovanda vesiculosa is a rootless aquatic plant. Seedlings develop a short protoroot; however, this fails to develop further and senesces. The plant consists of floating stems reaching a length of . The trap leaves grow in whorls of between 5 and 9 in close succession along the plant's central stem. The actual traps are held by petioles which have air sacs that aid in flotation. One end of the stem continually grows while the other end dies off. Growth is quite rapid ( per day in Japanese populations), so that in optimal conditions a new whorl is produced once or more each day.
Trap
The actual traps consist of two lobes which fold together to form a snap-trap similar to that of the Venus flytrap, except that it is smaller and located underwater. These traps, which are twisted so that the trap openings point outward, are lined on the inside by a fine coating of trigger hairs, snapping shut in response to contact with aquatic invertebrates and trapping them. The closing of this trap takes 10–20 milliseconds, making it one of the fastest examples of plant movement in the kingdom. This trapping is only possible in warm conditions of at least . Each trap is surrounded by between four and six long bristles that prevent triggering of traps by debris in the water.
Nutrient acquisition
A. vesiculosa is able to grow in nutrient-poor habitats not only due to its carnivory, but also due to its ability to re-utilize nutrients from senesced shoots, and its high affinity for mineral nutrients in water.
Reproduction
Flowers
The small, solitary white flowers of A. vesiculosa are supported above the water by short peduncles which arise from whorl axes. The flower only opens for a few hours, after which the structure is brought back beneath the water for seed production. The seeds are cryptocotylar: the cotyledons remain hidden within the seed coat and serve as an energy store for the seedling. Flowering, however, is rare in temperate regions and poorly successful in terms of fruit and seed development.
Divisions
Aldrovanda vesiculosa reproduces most often through vegetative reproduction. In favourable conditions, adult plants will produce an offshoot every , resulting in new plants as the tips continue to grow and the old ends die off and separate. Due to the rapid growth rate of this species, countless new plants can be produced in a short period of time in this fashion.
Turions
Winter-hardy Aldrovanda form turions as a frost survival strategy. At the onset of winter, the growth tip starts producing highly reduced non-carnivorous leaves on a severely shortened stem. This results in a tight bud of protective leaves which, being heavier and having released flotational gases, breaks off the mother plant and sinks to the water bottom, where temperatures are stable and warmer. Here it can withstand temperatures as low as . In the wild, Aldrovanda turions have been observed to have a relatively low rate of successful sinking. Those nutritious turions that fail to sink are then grazed by waterfowl or are killed by the onset of frost. In spring when water temperatures rise above , turions reduce their density and float to the top of the water, where they germinate and resume growth. Non-dormant turion-like organs can also form in response to summer drought.
Distribution
Aldrovanda vesiculosa is the second most widely distributed carnivorous plant species, only behind members of the genus Utricularia, native to Europe, Asia, Africa, and Australia. Aldrovanda is spread mainly through the movement of waterfowl: plants sticking to the feet of a bird are transported to the next aquatic destination on the bird's route. As a result, most Aldrovanda populations are located along avian migratory routes. Throughout the last century the species has become increasingly rare, listed as extinct in an increasingly large number of countries. In the 1970s, carnivorous plant hobbyists introduced this species to small backyard ponds in the United States in the states of New Jersey, Virginia, and the Catskills of New York, and they may be a potentially invasive species due to their effects on aquatic invertebrates.
Threats
Habitat degradation and modification
The waterwheel plant faces significant conservation threats related to habitat degradation and human-induced modifications. Residential and commercial development, along with agricultural and aquacultural activities, pose immediate risks to the species. The impacts of these activities on the plant's aquatic habitats are of particular concern.
Environmental changes
Across Europe, the species is confronted by several environmental challenges, as identified by the Commission of the European Union. Acidification, canalization, drainage, eutrophication, pollution, and various forms of habitat modification are highlighted as threats. These changes in the natural environment have the potential to disrupt the waterwheel plant's habitats and populations.
Potential illegal trade
While the extent and effects of illegal trade remain uncertain, it is believed that some illegal activities involving Aldrovanda vesiculosa may occur. This potential threat adds complexity to the conservation challenges faced by the species.
Habitat
A. vesiculosa prefers clean, shallow, warm, standing water with bright light, low nutrient levels, and a slightly acidic pH (around 6). It can be found floating amongst Juncus, reeds, and even rice.
The Waterwheel (Aldrovanda vesiculosa) thrives in a range of aquatic habitats, including small fens, peat-bog pools, billabongs, lakes, lagoons, and river deltas. It prefers oligo-mesotrophic and dystrophic systems with low nutrient levels. These plants are commonly found in shallow backwaters or the littoral zones of larger lakes, where they face less competition from other aquatic species and where water levels remain relatively stable throughout the growing season. The Waterwheel is highly intolerant of habitat degradation, and even minor changes in water chemistry can lead to local extinction.
Botanical history
Aldrovanda vesiculosa was first mentioned and illustrated in 1691 by Leonard Plukenet, based on collections made in India. He named the plant Lenticula palustris Indica. The modern botanical name originates from Gaetano Lorenzo Monti, who described Italian specimens in 1747 and named them Aldrovandia vesiculosa in honor of the Italian naturalist Ulisse Aldrovandi. When Carl Linnaeus published his Species Plantarum in 1753, the "i" was dropped from the name (an apparent copying error) to form the modern binomial.
Infraspecific taxa
Aldrovanda vesiculosa var. rubescens A.Cross and L.Adamec (2012)
Aldrovanda vesiculosa var. aquitanica Durieu ex Diels (1906) nom. illeg.
Aldrovanda vesiculosa var. australis Darwin (1876) nom. illeg.
Aldrovanda vesiculosa var. duriaei Caspary (1859) nom. illeg.
Aldrovanda vesiculosa var. verticillata (Roxb.) Darwin (1876) nom. illeg.
| Biology and health sciences | Caryophyllales | Plants |
873021 | https://en.wikipedia.org/wiki/Present | Present | The present is the period of time that is occurring now. The present is contrasted with the past, the period of time that has already occurred, and the future, the period of time that has yet to occur.
It is sometimes represented as a hyperplane in space-time, typically called "now", although modern physics demonstrates that such a hyperplane cannot be defined uniquely for observers in relative motion. The present may also be viewed as a duration.
Historiography
Contemporary history describes the historical timeframe immediately relevant to the present time and is a certain perspective of modern history.
Philosophy and religion
Philosophy of time
"The present" raises the question: "How is it that all sentient beings experience now at the same time?" There is no logical reason why this should be the case and no easy answer to the question.
In Buddhism
Buddhism and many of its associated paradigms emphasize the importance of living in the present moment—being fully aware of what is happening, and not dwelling on the past or worrying about the future. This does not mean that they encourage hedonism, but merely that constant focus on one's current position in space and time (rather than future considerations, or past reminiscence) will aid one in relieving suffering. They teach that those who live in the present moment are the happiest. A number of meditative techniques aim to help the practiser live in the present moment.
Christianity and eternity
Christianity views God as being outside of time and, from the divine perspective past, present and future are actualized in the now of eternity. This trans-temporal conception of God has been proposed as a solution to the problem of divine foreknowledge (i.e. how can God know what we will do in the future without us being determined to do it) since at least Boethius. Thomas Aquinas offers the metaphor of a watchman, representing God, standing on a height looking down on a valley to a road where past, present and future, represented by the individuals and their actions strung out along its length, are all visible simultaneously to God. Therefore, God's knowledge is not tied to any particular date.
Physical science
Special relativity
The original intent of the diagram on the right was to portray a 3-dimensional object having access to the past, present, and future in the present moment (4th dimension).
It follows from Albert Einstein's Special Theory of Relativity that there is no such thing as absolute simultaneity. When care is taken to operationalise "the present", it follows that the events that can be labeled as "simultaneous" with a given event, can not be in direct cause-effect relationship. Such collections of events are perceived differently by different observers. Instead, when focusing on "now" as the events perceived directly, not as a recollection or a speculation, for a given observer "now" takes the form of the observer's past light cone. The light cone of a given event is objectively defined as the collection of events in causal relationship to that event, but each event has a different associated light cone. One has to conclude that in relativistic models of physics there is no place for "the present" as an absolute element of reality, and only refers to things that are close to us. Einstein phrased this as: "People like us, who believe in physics, know that the distinction between past, present, and future is only a stubbornly persistent illusion".
Cosmology
In physical cosmology, the present time in the chronology of the universe is estimated at 13.8 billion years after the singularity determining the arrow of time.
In terms of the cosmic expansion history, it is in the dark-energy-dominated era, after the universe's matter content has become diluted enough for dark energy to dominate the total energy density. It is also in the universe's Stelliferous Era, after enough time for superclusters to have formed (at about 5 billion years), but before the accelerating expansion of the universe has removed the local supercluster beyond the cosmological horizon (at about 150 billion years).
Archaeology, geology, etc.
In radiocarbon dating, the "present" is defined as AD 1950.
Grammar
In English grammar, actions are classified according to one of the following twelve verb tenses: past (past, past continuous, past perfect, or past perfect continuous), present (present, present continuous, present perfect, or present perfect continuous), or future (future, future continuous, future perfect, or future perfect continuous). The present tense refers to things that are currently happening or are always the case. For example, in the sentence, "she walks home everyday," the verb "walks" is in the present tense because it refers to an action that is regularly occurring in the present circumstances.
Verbs in the present continuous tense indicate actions that are currently happening and will continue for a period of time. In the sentence, "she is walking home," the verb phrase "is walking" is in the present continuous tense because it refers to a current action that will continue until a certain endpoint (when "she" reaches home). Verbs in the present perfect tense indicate actions that started in the past and is completed at the time of speaking. For example, in the sentence, "She has walked home," the verb phrase "has walked" is in the present perfect tense because it describes an action that began in the past and is finished as of the current reference to the action. Finally, verbs in the present perfect continuous tense refer to actions that have been continuing up until the current time, thus combining the characteristics of both the continuous and perfect tenses. An example of a present perfect continuous verb phrase can be found in the sentence, "she has been walking this route for a week now," where "has been walking" indicates an action that was happening continuously in the past and continues to happen continuously in the present.
| Technology | Timekeeping | null |
874073 | https://en.wikipedia.org/wiki/Siren%20%28alarm%29 | Siren (alarm) | A siren is a loud noise-making device. There are two general types: mechanical and electronic. Civil defense sirens are mounted in fixed locations and used to warn of natural disasters or attacks. Sirens are used on emergency service vehicles such as ambulances, police cars, and fire engines.
Many fire sirens (used for summoning volunteer firefighters) serve double duty as tornado or civil defense sirens, alerting an entire community of impending danger. Most fire sirens are either mounted on the roof of a fire station or on a pole next to the fire station. Fire sirens can also be mounted on or near government buildings, on tall structures such as water towers, as well as in systems where several sirens are distributed around a town for better sound coverage. Most fire sirens are single tone and mechanically driven by electric motors with a rotor attached to the shaft. Some newer sirens are electronically driven speakers.
Fire sirens are often called fire whistles, fire alarms, or fire horns. Although there is no standard signaling of fire sirens, some utilize codes to inform firefighters of the location of the fire. Civil defense sirens also used as fire sirens often can produce an alternating "hi-lo" signal (similar to emergency vehicles in many European countries) as the fire signal, or attack (slow wail), typically 3x, as to not confuse the public with the standard civil defense signals of alert (steady tone) and fast wail (fast wavering tone). Fire sirens are often tested once a day at noon and are also called noon sirens or noon whistles.
The first emergency vehicles relied on a bell. In the 1970s, they switched to a duotone airhorn, which was itself overtaken in the 1980s by an electronic wail.
History
Some time before 1799, the siren was invented by the Scottish natural philosopher John Robison. Robison's sirens were used as musical instruments; specifically, they powered some of the pipes in an organ. Robison's siren consisted of a stopcock that opened and closed a pneumatic tube. The stopcock was apparently driven by the rotation of a wheel.
In 1819, an improved siren was developed and named by Baron Charles Cagniard de la Tour. De la Tour's siren consisted of two perforated disks that were mounted coaxially at the outlet of a pneumatic tube. One disk was stationary, while the other disk rotated. The rotating disk periodically interrupted the flow of air through the fixed disk, producing a tone. De la Tour's siren could produce sound under water, suggesting a link with the sirens of Greek mythology; hence the name he gave to the instrument.
Instead of disks, most modern mechanical sirens use two concentric cylinders, which have slots parallel to their length. The inner cylinder rotates while the outer one remains stationary. As air under pressure flows out of the slots of the inner cylinder and then escapes through the slots of the outer cylinder, the flow is periodically interrupted, creating a tone. The earliest such sirens were developed during 1877–1880 by James Douglass and George Slight (1859–1934) of Trinity House; the final version was first installed in 1887 at the Ailsa Craig lighthouse in Scotland's Firth of Clyde. When commercial electric power became available, sirens were no longer driven by external sources of compressed air, but by electric motors, which generated the necessary flow of air via a simple centrifugal fan, which was incorporated into the siren's inner cylinder.
To direct a siren's sound and to maximize its power output, a siren is often fitted with a horn, which transforms the high-pressure sound waves in the siren to lower-pressure sound waves in the open air.
The earliest way of summoning volunteer firemen to a fire was by ringing of a bell, either mounted atop the fire station, or in the belfry of a local church. As electricity became available, the first fire sirens were manufactured. In 1886 French electrical engineer Gustave Trouvé developed a siren to announce the silent arrival of his electric boats. Two early fire siren manufacturers were William A. Box Iron Works, who made the "Denver" sirens as early as 1905, and the Inter-State Machine Company (later the Sterling Siren Fire Alarm Company) who made the ubiquitous Model "M" electric siren, which was the first dual tone siren. The popularity of fire sirens took off by the 1920s, with many manufacturers including the Federal Electric Company and Decot Machine Works creating their own sirens. Since the 1970s, many communities have since deactivated their fire sirens as pagers became available for fire department use. Some sirens still remain as a backup to pager systems.
During the Second World War, the British civil defence used a network of sirens to alert the general population to the imminence of an air raid. A single tone denoted an "all clear". A series of tones denoted an air raid.
Types
Pneumatic
The pneumatic siren, which is a free aerophone, consists of a rotating disk with holes in it (called a chopper, siren disk or rotor), such that the material between the holes interrupts a flow of air from fixed holes on the outside of the unit (called a stator). As the holes in the rotating disk alternately prevent and allow air to flow it results in alternating compressed and rarefied air pressure, i.e. sound. Such sirens can consume large amounts of energy. To reduce the energy consumption without losing sound volume, some designs of pneumatic sirens are boosted by forcing compressed air from a tank that can be refilled by a low powered compressor through the siren disk.
In United States English language usage, vehicular pneumatic sirens are sometimes referred to as mechanical or coaster sirens, to differentiate them from electronic devices. Mechanical sirens driven by an electric motor are often called "electromechanical". One example is the Q2B siren sold by Federal Signal Corporation. Because of its high current draw (100 amps when power is applied) its application is normally limited to fire apparatus, though it has seen increasing use on type IV ambulances and rescue-squad vehicles. Its distinct tone of urgency, high sound pressure level (123 dB at 10 feet) and square sound waves account for its effectiveness.
In Germany and some other European countries, the pneumatic two-tone (hi-lo) siren consists of two sets of air horns, one high pitched and the other low pitched. An air compressor blows the air into one set of horns, and then it automatically switches to the other set. As this back and forth switching occurs, the sound changes tones. Its sound power varies, but could get as high as approximately 125 dB, depending on the compressor and the horns. Comparing with the mechanical sirens, it uses much less electricity but needs more maintenance.
In a pneumatic siren, the stator is the part which cuts off and reopens air as rotating blades of a chopper move past the port holes of the stator, generating sound. The pitch of the siren's sound is a function of the speed of the rotor and the number of holes in the stator. A siren with only one row of ports is called a single tone siren. A siren with two rows of ports is known as a dual tone siren.
By placing a second stator over the main stator and attaching a solenoid to it, one can repeatedly close and open all of the stator ports thus creating a tone called a pulse. If this is done while the siren is wailing (rather than sounding a steady tone) then it is called a pulse wail. By doing this separately over each row of ports on a dual tone siren, one can alternately sound each of the two tones back and forth, creating a tone known as Hi/Lo. If this is done while the siren is wailing, it is called a Hi/Lo wail. This equipment can also do pulse or pulse wail. The ports can be opened and closed to send Morse code. A siren which can do both pulse and Morse code is known as a code siren.
Electronic
Electronic sirens incorporate circuits such as oscillators, modulators, and amplifiers to synthesize a selected siren tone (wail, yelp, pierce/priority/phaser, hi-lo, scan, airhorn, manual, and a few more) which is played through external speakers. It is not unusual, especially in the case of modern fire engines, to see an emergency vehicle equipped with both types of sirens. Often, police sirens also use the interval of a tritone to help draw attention. The first electronic siren that mimicked the sound of a mechanical siren was invented in 1965 by Motorola employees Ronald H. Chapman and Charles W. Stephens.
Other types
Steam whistles were also used as a warning device if a supply of steam was present, such as a sawmill or factory. These were common before fire sirens became widely available, particularly in the former Soviet Union. Fire horns, large compressed air horns, also were and still are used as an alternative to a fire siren. Many fire horn systems were wired to fire pull boxes that were located around a town, and this would "blast out" a code in respect to that box's location. For example, pull box number 233, when pulled, would trigger the fire horn to sound two blasts, followed by a pause, followed by three blasts, followed by a pause, followed by three more blasts. In the days before telephones, this was the only way firefighters would know the location of a fire. The coded blasts were usually repeated several times. This technology was also applied to many steam whistles as well. Some fire sirens are fitted with brakes and dampers, enabling them to sound out codes as well. These units tended to be unreliable, and are now uncommon.
Physics of the sound
Mechanical sirens blow air through a slotted disk or rotor. The cyclic waves of air pressure are the physical form of sound. In many sirens, a centrifugal blower and rotor are integrated into a single piece of material, spun by an electric motor.
Electronic sirens are high efficiency loudspeakers, with specialized amplifiers and tone generation. They usually imitate the sounds of mechanical sirens in order to be recognizable as sirens.
To improve the efficiency of the siren, it uses a relatively low frequency, usually several hundred hertz. Lower frequency sound waves go around corners and through holes better.
Sirens often use horns to aim the pressure waves. This uses the siren's energy more efficiently by aiming it. Exponential horns achieve similar efficiencies with less material.
The frequency, i.e. the cycles per second of the sound of a mechanical siren is controlled by the speed of its rotor, and the number of openings. The wailing of a mechanical siren occurs as the rotor speeds and slows. Wailing usually identifies an attack or urgent emergency.
The characteristic timbre or musical quality of a mechanical siren is caused because it is a triangle wave, when graphed as pressure over time. As the openings widen, the emitted pressure increases. As they close, it decreases. So, the characteristic frequency distribution of the sound has harmonics at odd (1, 3, 5...) multiples of the fundamental. The power of the harmonics roll off in an inverse square to their frequency. Distant sirens sound more "mellow" or "warmer" because their harsh high frequencies are absorbed by nearby objects.
Two tone sirens are often designed to emit a minor third, musically considered a "sad" sound. To do this, they have two rotors with different numbers of openings. The upper tone is produced by a rotor with a count of openings divisible by six. The lower tone's rotor has a count of openings divisible by five. Unlike an organ, a mechanical siren's minor third is almost always physical, not tempered. To achieve tempered ratios in a mechanical siren, the rotors must either be geared, run by different motors, or have very large numbers of openings. Electronic sirens can easily produce a tempered minor third.
A mechanical siren that can alternate between its tones uses solenoids to move rotary shutters that cut off the air supply to one rotor, then the other. This is often used to identify a fire warning.
When testing, a frightening sound is not desirable. So, electronic sirens then usually emit musical tones: Westminster chimes is common. Mechanical sirens sometimes self-test by "growling", i.e. operating at low speeds.
In music
Sirens are also used as musical instruments. They have been prominently featured in works by avant-garde and contemporary classical composers. Examples include Edgard Varèse's compositions Amériques (1918–21, rev. 1927), Hyperprism (1924), and Ionisation (1931); Arseny Avraamov's Symphony of Factory Sirens (1922); George Antheil's Ballet Mécanique (1926); Dimitri Shostakovich's Symphony No. 2 (1927), and Henry Fillmore's "The Klaxon: March of the Automobiles" (1929), which features a klaxophone.
In popular music, sirens have been used in The Chemical Brothers' "Song to the Siren" (1992) and in a CBS News 60 Minutes segment played by percussionist Evelyn Glennie. A variation of a siren, played on a keyboard, are the opening notes of the REO Speedwagon song "Ridin' the Storm Out". Some heavy metal bands also use air raid type siren intros at the beginning of their shows. The opening measure of Money City Maniacs 1998 by Canadian band Sloan uses multiple sirens overlapped.
Vehicle-mounted
Approvals or certifications
Governments may have standards for vehicle-mounted sirens. For example, in California, sirens are designated Class A or Class B. A Class A siren is loud enough that it can be mounted nearly anywhere on a vehicle. Class B sirens are not as loud and must be mounted on a plane parallel to the level roadway and parallel to the direction the vehicle travels when driving in a straight line.
Sirens must also be approved by local agencies, in some cases. For example, the California Highway Patrol approves specific models for use on emergency vehicles in the state. The approval is important because it ensures the devices perform adequately. Moreover, using unapproved devices could be a factor in determining fault if a collision occurs.
The SAE International Emergency Warning Lights and Devices committee oversees the SAE emergency vehicle lighting practices and the siren practice, J1849. This practice was updated through cooperation between the SAE and the National Institute of Standards and Technology. Though this version remains quite similar to the California Title 13 standard for sound output at various angles, this updated practice enables an acoustic laboratory to test a dual speaker siren system for compliant sound output.
Best practices
Siren speakers, or mechanical sirens, should always be mounted ahead of the passenger compartment. This reduces the noise for occupants and makes two-way radio and mobile telephone audio more intelligible during siren use. It also puts the sound where it will be useful. A 2007 study found passenger compartment sound levels could exceed 90dB(A).
Research has shown that sirens mounted behind the engine grille or under the wheel arches produces less unwanted noise inside the passenger cabin and to the side and rear of the vehicle while maintaining noise levels to give adequate warnings. The inclusion of broadband sound to sirens has the ability to increase localisation of sirens, as in a directional siren, as a spread of frequencies makes use of the three ways the brain detects a direction of a sound: Interaural level difference, interaural time difference and head-related transfer function.
The worst installations are those where the siren sound is emitted above and slightly behind the vehicle occupants such as cases where a light-bar mounted speaker is used on a sedan or pickup. Vehicles with concealed sirens also tend to have high noise levels inside. In some cases, concealed or poor installations produce noise levels which can permanently damage vehicle occupants' hearing.
Electric-motor-driven mechanical sirens may draw 50 to 200 amperes at 12 volts (DC) when spinning up to operating speed. Appropriate wiring and transient protection for engine control computers is a necessary part of an installation. Wiring should be similar in size to the wiring to the vehicle engine starter motor. Mechanical vehicle mounted devices usually have an electric brake, a solenoid that presses a friction pad against the siren rotor. When an emergency vehicle arrives on-scene or is cancelled en route, the operator can rapidly stop the siren.
Multi-speaker electronic sirens often are alleged to have dead spots at certain angles to the vehicle's direction of travel. These are caused by phase differences. The sound coming from the speaker array can phase cancel in some situations. This phase cancellation occurs at single frequencies, based upon the spacing of the speakers. These phase differences also account for increases, based upon the frequency and the speaker spacing. However, sirens are designed to sweep the frequency of their sound output, typically, no less than one octave. This sweeping minimizes the effects of phase cancellation. The result is that the average sound output from a dual speaker siren system is 3 dB greater than a single speaker system.
| Technology | Mechanisms | null |
20077968 | https://en.wikipedia.org/wiki/Hyssopus%20officinalis | Hyssopus officinalis | Hyssopus officinalis or hyssop is a shrub in the Lamiaceae or mint family native to Southern Europe, the Middle East, and the region surrounding the Caspian Sea. Due to its purported properties as an antiseptic, cough reliever, and expectorant, it has been used in traditional herbal medicine.
Description
Hyssop is a brightly coloured shrub or subshrub that ranges from in height. The stem is woody at the base, from which grow a number of upright branches. Its leaves are lanceolate, dark green, and from long.
During the summer, hyssop produces pink, blue (ssp. aristadus), or, more rarely, white fragrant (ssp. f. albus) flowers. These give rise to small oblong tetra-achenes.
History
A plant called hyssop has been in use since classical antiquity. Its name is a direct adaptation from the Greek ὕσσωπος (). The Hebrew word אזוב (ezov, esov, or esob) and the Greek word ὕσσωπος probably share a common (but unknown) origin. The name hyssop appears as a translation of ezov in some translations of the Bible, notably in : "Purge me with hyssop, and I shall be clean", but researchers have suggested that the Biblical accounts refer not to the plant currently known as hyssop but rather to one of a number of different herbs, including Origanum syriacum (Syrian oregano, commonly referred to as "bible hyssop"). mentions that 'ezov' was a small plant and some scholars believe it was a wall plant. It was burned with the red heifer () and used for purification of lepers (, ; ), and at Passover it was used to sprinkle the blood of the sacrificial lamb on the doorposts (). A sponge attached to a hyssop branch was used to give Jesus on the cross a drink of vinegar.
Suggestions abound for the modern day correlation of biblical hyssop ranging from a wall plant like moss or fern, to widely used culinary herbs like thyme, rosemary or marjoram. Another suggestion is the caper plant which is known to grow in the rocky soils of the region and along walls.
Hyssop was also used for purgation (religious purification) in Egypt, where, according to Chaeremon the Stoic, the priests used to eat it with bread in order to purify this type of food and make it suitable for their austere diet.
Cultivation
Hyssop is resistant to drought, and tolerant of chalky, sandy soils. It thrives in full sun and warm climates.
Cultivars include 'Blue Flower'.
Harvest
Under optimal weather conditions, herb hyssop is harvested twice yearly, once at the end of spring and once more at the beginning of autumn. The plants are preferably harvested when flowering in order to collect the flowering tips.
Once the stalks are cut, they are collected and dried either stacked on pallets to allow for draining or hung to dry. The actual drying process takes place in a cool, dry, well-ventilated area, where the materials are mixed several times to ensure even drying. Drying herbs are kept from exposure to the sun to prevent discoloration and oxidation. The drying process takes approximately six days in its entirety. Once dried, the leaves are removed and both components, leaves and flowers, are chopped finely. The final dried product weighs a third of the initial fresh weight and can be stored for up to 18 months.
Essential oil
The essential oil includes the chemicals thujone and phenol, which give it antiseptic properties. Its high concentrations of thujone and chemicals that stimulate the central nervous system, including pinocamphone and cineole, can provoke epileptic reactions. The oil of hyssop can cause seizures and even low doses (2–3 drops) can cause convulsions in children.
Uses
Culinary
The fresh herb is commonly used in cooking. Za'atar is a famous Middle Eastern herbal mixture, some versions of which include dried hyssop leaves.
Essence of hyssop can be obtained by steaming, and is used in cooking to a lesser extent.
The plant is commonly used by beekeepers to produce nectar from which western honey bees make a rich and aromatic honey.
Herb hyssop leaves are used as an aromatic condiment. The leaves have a lightly bitter taste due to its tannins, and an intense minty aroma. Due to its intensity, it is used moderately in cooking. The herb is also used to flavor liqueur, and is part of the official formulation of Chartreuse. It is also a key ingredient in many formulations of absinthe, where it is the main source of the green colour.
Herbal medicine
In herbal medicine hyssop is believed to have soothing, expectorant, and cough suppressant properties. Hyssop has been used for centuries in traditional medicine in order to increase circulation and to treat multiple conditions such as coughing and sore throat. Hyssop can stimulate the gastrointestinal system.
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| Biology and health sciences | Herbs and spices | Plants |
8543506 | https://en.wikipedia.org/wiki/Cave%20hyena | Cave hyena | Cave hyena (Crocuta (crocuta) spelaea and Crocuta (crocuta) ultima) are extinct species or subspecies of hyena known from Eurasia, which ranged from Western Europe to eastern Asia and Siberia during the Pleistocene epoch. It is well represented in many European caves, primarily dating to the Last Glacial Period. It was an apex predator that preyed on large mammals (primarily large ungulates, such as wild horse and aurochs), and was responsible for the accumulation of hundreds of large Pleistocene mammal bones in areas including horizontal caves, sinkholes, mud pits, and muddy areas along rivers.
Often treated as subspecies or populations of the living African spotted hyena (Crocuta crocuta) to which they were closely related and heavily resembled, genetic evidence from the nuclear genome suggests that Eurasian Crocuta populations (including the west Eurasian Crocuta crocuta spelaea and Asian Crocuta crocuta ultima) were highly genetically divergent from African populations (having estimated to have split over 1 million years ago), with evidence suggesting limited interbreeding between Eurasian cave and African spotted hyenas. Some authors have suggested that the two subspecies should be raised to species level as Crocuta spelaea and Crocuta ultima.
Cave hyenas coexisted in Europe alongside both Neanderthals and modern humans. Evidence suggests that they consumed the remains of Neanderthals at least on occasion, with cave hyenas also being recorded in cave art.
The cause of the cave hyena's extinction is not fully understood, though it could have been due to a combination of factors, including human activity, diminished quantities of prey animals, and climate change.
Description
The size of cave hyenas varied depending on environment, with populations inhabiting colder climates having a larger body size than those inhabiting more temperate climates, an example of Bergmann's rule. A 2017 study estimated that on average cave hyenas weighed approximately , around 60% heavier than living spotted hyenas. In comparison to living spotted hyenas, some of the bones of the limbs are more robust (proportionally thicker and shorter), with the ulna being more curved. In the skull, the first and second upper premolars contact each other, while in living spotted hyenas they are separated by a diastema (gap), though collectively the differences between the skeletal anatomy of cave hyenas and living spotted hyenas are "relatively minor". The endocast (brain cavity) also shows differences between from that of the spotted hyena, with the living spotted hyena having a better developed front part of the brain, which may suggest differences in behaviour. Evidence from cave art suggests that cave hyenas had a similar physical appearance to living spotted hyenas, with a spotted pelt.
Ecology
While "cave hyenas" did use caves, they were not confined to them, being present where caves were absent, and where present only using caves intermittently, with some examples of open-air dens having being found in the fossil record, such as at Bottrop in Germany.
The cave hyena's diet probably differed little from contemporary African spotted hyenas, and like living spotted hyenas, cave hyenas probably lived in groups (which in living spotted hyenas are called "clans") and were active predators rather than purely scavengers (with hunting being predominant over scavenging in the living spotted hyena). The bones of cave hyena prey were often cracked open/crushed in order to feed on the interior marrow, as is done by living spotted hyenas.
The diet of the cave hyena is thought to have primarily consisted of large ungulates like wild horse, aurochs, steppe bison, Irish elk/giant deer, wild boar, red deer and reindeer, with larger herbivores like the woolly rhinoceros and woolly mammoth probably being scavenged after death (or alternatively generally only targeted when weak), with their young perhaps sometimes being targeted for active hunting.
Wild horse are common in Late Pleistocene European cave hyena dens, implying that they were frequent prey, as zebras are for modern African spotted hyenas. In karst cave sites in the Czech Republic which accumulated during the Last Glacial Period, wild horses as well as woolly rhinoceros are the most common remains in cave hyena dens (though this to a significant degree probably reflects the durability of woolly rhinoceros bones making them able to be identified), with other remains including steppe bison, red deer, reindeer, European wild ass, chamois, alpine ibex, and cave bears (who may have been scavenged after dying in the caves). At the Fouvent-Saint-Andoche karst in France, a similar assemblage is found, albeit also including juvenile woolly mammoths, reflecting the similar cold conditions in both regions during the Last Glacial Period. At Perspektywiczna Cave in Poland, the assemblage is dominated by reindeer, with other remains present including bovids, equines, woolly rhinoceros and woolly mammoth.
At the Grotta Paglicci den in southern Italy, the assemblage is dominated by European fallow deer, red deer, roe deer and aurochs. At San Teodoro cave in northern Sicily which is a well known hyena den, remains of herbivores likely accumulated in the cave by hyenas include those of aurochs, steppe bison (with aurochs seeming to predominate over bison), European wild ass, wild boar, red deer, and the endemic dwarf elephant Palaeoloxodon mnaidriensis. At Kirkdale Cave in Yorkshire, northern England which dates to the Last Interglacial when Europe had a temperate climate similar to modern times, the assemblage includes juvenile straight-tusked elephant, Irish elk, red deer, European fallow deer, bison, and the narrow-nosed rhinoceros. At the Manot Cave site in Israel, the bone assemblage accumulated by cave hyenas is predominantly Persian fallow deer, as well as to a lesser extent goats, aurochs and equines. At Wezmeh Cave in the Zagros Mountains of western Iran, dating to the Last Glacial Period, remains include mouflon, wild goat, red deer, aurochs, wild boar, gazelles, and the narrow-nosed rhinoceros. Cave hyenas also appear to have extensively engaged in cannibalism, demonstrated as numerous sites.
Cave hyenas likely came into conflict with cave lions (which despite their name, probably only rarely if ever used caves) over carcasses, with remains of cave lions found in European cave deposits possibly being the result of being brought into the cave by cave hyenas. A significant number of these cave lion remains do not bear any evidence of consumption, which may suggest that like living spotted hyenas, they did not generally consume the remains of lions after killing them.
History of discovery and classification
In the 1770s, Zoolithen Cave in Germany gained scientific attention due to the remains of fossil animals found within. In 1822 a skull from this cave was used to describe the species Hyena spelaea by German paleontologist Georg August Goldfuss.
Although the first full account of the cave hyena was given by Georges Cuvier in 1812, skeletal fragments of the cave hyena have been described in scientific literature since the 18th century, though they were frequently misidentified. The first recorded mention of the cave hyena in literature occurs in Kundmann's 1737 tome Rariora Naturæ et Artis, where the author misidentified a hyena's mandibular ramus as that of a calf. In 1774, Esper erroneously described hyena teeth discovered in Gailenreuth as those of a lion, and in 1784, Collini described a cave hyena skull as that of a seal. The past presence of hyenas in Great Britain was revealed after William Buckland's examination of the contents of Kirkdale Cave, which was discovered to have once been the location of several hyena den sites. Buckland's findings were followed by further discoveries by Clift and Whidbey in Oreston, Plymouth.
In his own 1812 account, Cuvier mentioned a number of European localities where cave hyena remains were found, and considered it a different species from the spotted hyena on account of its superior size. He elaborated his view in his Ossemens Fossiles (1823), noting how the cave hyena's digital extremities were shorter and thicker than those of the spotted hyena. His views were largely accepted throughout the first half of the 19th century, finding support in de Blainville and Richard Owen among others. Further justifications in separating the two animals included differences in the tubercular portion of the lower carnassial. Boyd Dawkins, writing in 1865, was the first to definitely cast doubt over the separation of the spotted and cave hyena, stating that the aforementioned tooth characteristics were consistent with mere individual variation. Writing again in 1877, he further stated after comparing the two animals' skulls that there are no characters of specific value.
Genetics
Analysis of the mitochondrial genomes of Eurasian Crocuta specimens shows no clear separation from African lineages. However, an analysis of full nuclear genomes of both European and East Asian cave hyenas published in 2020 suggests that African and Eurasian Crocuta populations were largely separate, having estimated to have diverged from each other around 2.5 million years ago, closely corresponding to the age of the earliest Crocuta specimens in Eurasia, which are around 2 million years old from China. The nuclear genome results also suggested that the European and East Asian populations (often assigned to the separate subspecies C. crocuta ultima) were strongly genetically divergent from each other, but were more closely related to each other overall than to African Crocuta populations. Analysis of the nuclear genome suggests that there had been interbreeding between these populations for some time after the split, which likely explains the discordance between the nuclear and mitochondrial genome results, with the mitochondrial genomes of African and European Crocuta more closely related to each other than to East Asian Crocuta, suggesting gene flow between the two groups after the split between the East Asian and European populations. Some authors have suggested that the two subspecies should be raised to species level as Crocuta spelaea and Crocuta ultima. A 2024 study of a cave hyena genome from Sicily found that as with the 2020 study, there was strong genetic separation between Eurasian cave and African spotted hyenas, but unlike the 2020 study, there was no robust support for a basal split between East Asian/Siberian and European cave hyenas, with the Sicilian cave hyena found to be the earliest diverging cave hyena lineage, with less interbreeding with African hyenas than other European cave hyenas.
Distribution
Crocuta first appeared outside of Africa in Asia during the Early Pleistocene around 2 million years ago, before arriving in Europe at the beginning of the Middle Pleistocene around 800,000 years ago, around the time of the extinction of the "giant hyena" Pachycrocuta brevirostris in the region. Crocuta was widely distributed across northern Eurasia during the Middle-Late Pleistocene, spanning from the Iberian Peninsula, Britain and Ireland in the West, across southern Siberia, Mongolia and northern China to the Pacific Coast of the Russian Far East. C. c. ultima at times ranged as far southeast as Guangxi and Taiwan in southern China, as well as Thailand and Laos in Southeast Asia, while C. c. spelaea ranged into the Middle East, as far south as the Judaean Desert and as far east as the Zagros Mountains of western Iran. The northernmost records are from the banks of the Vilyuy River in Northeast Siberia, with indirect evidence of feeding on woolly rhinoceros carcasses suggesting cave hyenas may have reached the far northeast of Siberia near the Arctic circle.
Relationships with hominids
Interactions
Kills partially processed by Neanderthals and then by cave hyenas indicate that hyenas would occasionally steal Neanderthal kills; and cave hyenas and Neanderthals competed for cave sites. Many caves show alternating occupations by hyenas and Neanderthals. There is fossil evidence of humans in Middle Pleistocene Europe butchering and presumably consuming hyenas. At a number of cave hyena den sites, the remains of Neanderthals have been found showing evidence of having been gnawed on by the hyenas, this may be the result of cave hyenas scavenging Neanderthal burials, though some of these remains may also be the result of cave hyenas attacking and killing Neanderthals.
In rock art
The cave hyena is depicted in a few examples of Upper Palaeolithic rock art in France. A painting from the Chauvet Cave depicts a hyena outlined and represented in profile, with two legs, with its head and front part with well distinguishable spotted coloration pattern. Because of the specimen's steeped profile, it is thought that the painting was originally meant to represent a cave bear, but was modified as a hyena. In Lascaux, a red and black rock painting of a hyena is present in the part of the cave known as the Diverticule axial, and is depicted in profile, with four limbs, showing an animal with a steep back. The body and the long neck have spots, including the flanks. An image on a cave in Ariège shows an incompletely outlined and deeply engraved figure, representing a part of an elongated neck, smoothly passing into part of the animal's forelimb on the proximal side. Its head is in profile, with a possibly re-engraved muzzle. The ear is typical of the spotted hyena, as it is rounded. An image in the Le Gabillou Cave in Dordogne shows a deeply engraved zoomorphic figure with a head in frontal view and an elongated neck with part of the forelimb in profile. It has large round eyes and short, rounded ears which are set far from each other. It has a broad, line-like mouth that evokes a smile. Though originally thought to represent a composite or zoomorphic hybrid, it is probable it is a spotted hyena based on its broad muzzle and long neck. The relative scarcity of hyena depictions in Paleolithic rock art has been theorised to be due to the animal's lower rank in the animal worship hierarchy; the cave hyena's appearance was likely unappealing to Ice Age hunters, and it was not sought after as prey. Also, it was not a serious rival like the cave lion or bear, and it lacked the impressiveness of the mammoth or woolly rhino.
Extinction
A 2014 study concluded that the youngest well-dated remains of cave hyenas in Europe date to around 31,000 years ago. Subsequent studies have suggested that cave hyenas may have persisted later in the Iberian Peninsula based on the radiocarbon dating of coprolites attributed to them, possibly as late as 7,000 years ago in the southern Iberian Peninsula, but it is suggested that the dates should be considered with caution due to potential issues with contamination. A 2021 study found the youngest specimens in East Asia date to around 20,000 years ago. Potential causal factors for extinction include decreasing temperatures, competition with other carnivores, including humans for food and living space, and decreased prey abundance. Evidence suggests that climate change alone cannot account for the cave hyena's extinction in Europe and that other factors, such as human activity and decreasing prey abundance, are necessary to explain it.
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| Biology and health sciences | Other carnivora | Animals |
519679 | https://en.wikipedia.org/wiki/Diol | Diol | A diol is a chemical compound containing two hydroxyl groups ( groups). An aliphatic diol may also be called a glycol. This pairing of functional groups is pervasive, and many subcategories have been identified. They are used as protecting groups of carbonyl groups, making them essential in synthesis of organic chemistry.
The most common industrial diol is ethylene glycol. Examples of diols in which the hydroxyl functional groups are more widely separated include 1,4-butanediol and propylene-1,3-diol, or beta propylene glycol, .
Synthesis of classes of diols
Geminal diols
A geminal diol has two hydroxyl groups bonded to the same atom. These species arise by hydration of the carbonyl compounds. The hydration is usually unfavorable, but a notable exception is formaldehyde which, in water, exists in equilibrium with methanediol H2C(OH)2. Another example is (F3C)2C(OH)2, the hydrated form of hexafluoroacetone. Many gem-diols undergo further condensation to give dimeric and oligomeric derivatives. This reaction applies to glyoxal and related aldehydes.
Vicinal diols
In a vicinal diol, the two hydroxyl groups occupy vicinal positions, that is, they are attached to adjacent atoms. These compounds are called glycols (though the term can be used more widely). Examples include ethane-1,2-diol or ethylene glycol HO−(CH2)2−OH, a common ingredient of antifreeze products. Another example is propane-1,2-diol, or alpha propylene glycol, HO−CH2−CH(OH)−CH3, used in the food and medicine industry, as well as a relatively non-poisonous antifreeze product.
On commercial scales, the main route to vicinal diols is the hydrolysis of epoxides. The epoxides are prepared by epoxidation of the alkene. An example in the synthesis of trans-cyclohexanediol or by microreactor:
For academic research and pharmaceutical areas, vicinal diols are often produced from the oxidation of alkenes, usually with dilute acidic potassium permanganate or Osmium tetroxide. Osmium tetroxide can similarly be used to oxidize alkenes to vicinal diols. The chemical reaction called Sharpless asymmetric dihydroxylation can be used to produce chiral diols from alkenes using an osmate reagent and a chiral catalyst. Another method is the Woodward cis-hydroxylation (cis diol) and the related Prévost reaction (anti diol), which both use iodine and the silver salt of a carboxylic acid.
Other routes to vic-diols are the hydrogenation of acyloins and the pinacol coupling reaction.
1,3-Diols
1,3-Diols are often prepared industrially by aldol condensation of ketones with formaldehyde. You can use many different starting materials to produce syn- or anti-1,3-diols. The resulting carbonyl is reduced using the Cannizzaro reaction or by catalytic hydrogenation:
RC(O)CH3 + CH2O → RC(O)CH2CH2OH
RC(O)CH2CH2OH + H2 → RCH(OH)CH2CH2OH
2,2-Disubstituted propane-1,3-diols are prepared in this way. Examples include 2-methyl-2-propyl-1,3-propanediol and neopentyl glycol.
1,3-Diols can be prepared by hydration of α,β-unsaturated ketones and aldehydes. The resulting keto-alcohol is hydrogenated. Another route involves the hydroformylation of epoxides followed by hydrogenation of the aldehyde. This method has been used for 1,3-propanediol from ethylene oxide.
More specialized routes to 1,3-diols involves the reaction between an alkene and formaldehyde, the Prins reaction. 1,3-diols can be produced diastereoselectively from the corresponding β-hydroxy ketones using the Evans–Saksena, Narasaka–Prasad or Evans–Tishchenko reduction protocols.
1,3-Diols are described as syn or anti depending on the relative stereochemistries of the carbon atoms bearing the hydroxyl functional groups. Zincophorin is a natural product that contains both syn and anti 1,3-diols.
1,4-, 1,5-, and longer diols
Diols where the hydroxyl groups are separated by several carbon centers are generally prepared by hydrogenation of diesters of the corresponding dicarboxylic acids:
(CH2)n(CO2R)2 + 4 H2 → (CH2)n(CH2OH)2 + 2 H2O + 2 ROH
1,4-butanediol, 1,5-pentanediol, 1,6-hexanediol, and are important precursors to polyurethanes.
Reactions
From the industrial perspective, the dominant reactions of the diols is in the production of polyurethanes and alkyd resins.
General diols
Diols react as alcohols, by esterification and ether formation.
Diols such as ethylene glycol are used as co-monomers in polymerization reactions forming polymers including some polyesters and polyurethanes. A different monomer with two identical functional groups, such as a dioyl dichloride or dioic acid is required to continue the process of polymerization through repeated esterification processes.
A diol can be converted to cyclic ether by using an acid catalyst, this is diol cyclization. Firstly, it involves protonation of the hydroxyl group. Then, followed by intramolecular nucleophilic substitution, the second hydroxyl group attacks the electron deficient carbon. Provided that there are enough carbon atoms that the angle strain is not too much, a cyclic ether can be formed.
1,2-diols and 1,3-diols can be protected using a protecting group. Protecting groups are used so that the functional group does not react to future reactions. Benzylidene groups are used to protect 1,3-diols. There are extremely useful in biochemistry as shown below of a carbohydrate derivative being protected.
Diols can also be used to protect carbonyl groups. They are commonly used and are quite efficient at synthesizing cyclic acetals. These protect the carbonyl groups from reacting from any further synthesis until it is necessary to remove them. The reaction below depicts a diol being used to protect a carbonyl using zirconium tetrachloride.
Diols can also be converted to lactones employing the Fétizon oxidation reaction.
Vicinal diols
In glycol cleavage, the C−C bond in a vicinal diol is cleaved with formation of ketone or aldehyde functional groups. See Diol oxidation.
Geminal diols
In general, organic geminal diols readily dehydrate to form a carbonyl group.
| Physical sciences | Alcohols | Chemistry |
519796 | https://en.wikipedia.org/wiki/Iodide | Iodide | An iodide ion is the ion I−. Compounds with iodine in formal oxidation state −1 are called iodides. In everyday life, iodide is most commonly encountered as a component of iodized salt, which many governments mandate. Worldwide, iodine deficiency affects two billion people and is the leading preventable cause of intellectual disability.
Structure and characteristics of inorganic iodides
Iodide is one of the largest monatomic anions. It is assigned a radius of around 206 picometers. For comparison, the lighter halides are considerably smaller: bromide (196 pm), chloride (181 pm), and fluoride (133 pm). In part because of its size, iodide forms relatively weak bonds with most elements.
Most iodide salts are soluble in water, but often less so than the related chlorides and bromides. Iodide, being large, is less hydrophilic compared to the smaller anions. One consequence of this is that sodium iodide is highly soluble in acetone, whereas sodium chloride is not. The low solubility of silver iodide and lead iodide reflects the covalent character of these metal iodides. A test for the presence of iodide ions is the formation of yellow precipitates of these compounds upon treatment of a solution of silver nitrate or lead(II) nitrate.
Aqueous solutions of iodide salts dissolve iodine better than pure water. This effect is due to the formation of the triiodide ion, which is brown:
I− + I2 ⇌
Redox, including antioxidant properties
Iodide salts are mild reducing agents and many react with oxygen to give iodine. A reducing agent is a chemical term for an antioxidant. Its antioxidant properties can be expressed quantitatively as a redox potential :
2I− ⇌ I2 + E° = 0.54 volts (versus SHE)
Because iodide is easily oxidized, some enzymes readily convert it into electrophilic iodinating agents, as required for the biosynthesis of myriad iodide-containing natural products. Iodide can function as an antioxidant reducing species that can destroy ozone and reactive oxygen species such as hydrogen peroxide:
2 I− + peroxidase + H2O2 + tyrosine, histidine, lipid, etc. → iodo-compounds + H2O + 2 e− (antioxidants).
Representative iodides
Natural occurrence
Iodargyrite—natural, crystalline silver iodide—is the most common iodide mineral currently known. Iodide anions may sometimes also be found combined with mercury, copper and lead, but minerals with such compositions are even more scarce.
Other oxoanions
Iodine can assume oxidation states of −1, +1, +3, +5, or +7. A number of neutral iodine oxides are also known.
| Physical sciences | Halide salts | Chemistry |
519841 | https://en.wikipedia.org/wiki/Nucleophilic%20addition | Nucleophilic addition | In organic chemistry, a nucleophilic addition (AN) reaction is an addition reaction where a chemical compound with an electrophilic double or triple bond reacts with a nucleophile, such that the double or triple bond is broken. Nucleophilic additions differ from electrophilic additions in that the former reactions involve the group to which atoms are added accepting electron pairs, whereas the latter reactions involve the group donating electron pairs.
Addition to carbon–heteroatom double bonds
Nucleophilic addition reactions of nucleophiles with electrophilic double or triple bond (π bonds) create a new carbon center with two additional single, or σ, bonds. Addition of a nucleophile to carbon–heteroatom double or triple bonds such as >C=O or -C≡N show great variety. These types of bonds are polar (have a large difference in electronegativity between the two atoms); consequently, their carbon atoms carries a partial positive charge. This makes the molecule an electrophile, and the carbon atom the electrophilic center; this atom is the primary target for the nucleophile. Chemists have developed a geometric system to describe the approach of the nucleophile to the electrophilic center, using two angles, the Bürgi–Dunitz and the Flippin–Lodge angles after scientists that first studied and described them.
This type of reaction is also called a 1,2-nucleophilic addition. The stereochemistry of this type of nucleophilic attack is not an issue, when both alkyl substituents are dissimilar and there are not any other controlling issues such as chelation with a Lewis acid, the reaction product is a racemate. Addition reactions of this type are numerous. When the addition reaction is accompanied by an elimination the reaction is a type of substitution or an addition-elimination reaction.
Addition to carbonyl groups
With a carbonyl compound as an electrophile, the nucleophile can be:
water in hydration to a geminal diol (hydrate)
an alcohol in acetalisation to an acetal
a hydride in reduction to an alcohol
an amine with formaldehyde and a carbonyl compound in the Mannich reaction
an enolate ion in an aldol reaction or Baylis–Hillman reaction
an organometallic nucleophile in the Grignard reaction or the related Barbier reaction or a Reformatskii reaction
ylides such as a Wittig reagent or the Corey–Chaykovsky reagent or α-silyl carbanions in the Peterson olefination
a phosphonate carbanion in the Horner–Wadsworth–Emmons reaction
a pyridine zwitterion in the Hammick reaction
an acetylide in alkynylation reactions.
a cyanide ion in cyanohydrin reactions
In many nucleophilic reactions, addition to the carbonyl group is very important. In some cases, the C=O double bond is reduced to a C-O single bond when the nucleophile bonds with carbon. For example, in the cyanohydrin reaction a cyanide ion forms a C-C bond by breaking the carbonyl's double bond to form a cyanohydrin.
Addition to nitriles
With nitrile electrophiles, nucleophilic addition take place by:
hydrolysis of a nitrile to form an amide or a carboxylic acid
organozinc nucleophiles in the Blaise reaction
alcohols in the Pinner reaction.
the (same) nitrile α-carbon in the Thorpe reaction. The intramolecular version is called the Thorpe–Ziegler reaction.
Grignard reagents to form imines. The route affords ketones following hydrolysis or primary amines following imine reduction.
Addition to carbon–carbon double bonds
When a nucleophile X− adds to an alkene, the driving force is the transfer of negative charge from X to the electron-poor unsaturated -C=C- system. This occurs through the formation of a covalent bond between X and one carbon atom, concomitant with the transfer of electron density from the pi bond onto the other carbon atom (step 1). During a telescoped second reaction or workup (step 2), the resulting negatively charged carbanion combines with an electrophilic Y to form the second covalent bond.
Unsubstituted and unstrained alkenes are typically insufficiently polar to admit nucleophilic addition, but a few exceptions are known.
The strain energy in fullerenes weakens their double-bonds; addition thereto is the Bingel reaction.
Bonds adjacent to an electron-withdrawing substituent (e.g. a carbonyl group, nitrile, or fluoride) readily admit nucleophilic addition. In this process, conjugate addition, the nucleophile X adds β to the substituent, because then said substituent inductively stabilizes the product's negative charge. Aromatic substituents, although typically electrophilic, can also sometimes stabilize negative charge; for example, styrene reacts in toluene with sodium to give 1,3-diphenylpropane:
| Physical sciences | Organic reactions | Chemistry |
519887 | https://en.wikipedia.org/wiki/Fiber%20crop | Fiber crop | Fiber crops are field crops grown for their fibers, which are traditionally used to make paper, cloth, or rope.
Fiber crops are characterized by having a large concentration of cellulose, which is what gives them their strength. The fibers may be chemically modified, like in viscose (used to make rayon and cellophane). In recent years, materials scientists have begun exploring further use of these fibers in composite materials. Due to cellulose being the main factor of a plant fiber's strength, this is what scientists are looking to manipulate to create different types of fibers.
Fiber crops are generally harvestable after a single growing season, as distinct from trees, which are typically grown for many years before being harvested for such materials as wood pulp fiber or lacebark. In specific circumstances, fiber crops can be superior to wood pulp fiber in terms of technical performance, environmental impact or cost.
There are a number of issues regarding the use of fiber crops to make pulp. One of these is seasonal availability. While trees can be harvested continuously, many field crops are harvested once during the year and must be stored such that the crop doesn't rot over a period of many months. Considering that many pulp mills require several thousand tonnes of fiber source per day, storage of the fiber source can be a major issue.
Botanically, the fibers harvested from many of these plants are bast fibers; the fibers come from the phloem tissue of the plant. The other fiber crop fibers are hard/leaf fibers (come from the entirety of plant vascular bundles) and surface fibers (come from plant epidermal tissue).
Fiber sources
To have a source of fiber to utilize in production, the fiber first must be extracted from the plant. This is done in different ways depending on the fiber classification. Bast fibers are harvested through retting which is where microbes are utilized to remove soft tissues from the plant and only the useful fibrous material remains. Hard fibers are harvested mainly through decortication which is where the non-fibrous tissues are removed by hand or machine. Lastly, surface fibers are harvested through ginning which is where a machine removes the fibers from other plant material.
Paper
Before the industrialisation of paper production the most common fiber source was recycled fibers from used textiles, called rags. The rags were from ramie, hemp, linen and cotton. A process for removing printing inks from recycled paper was invented by German jurist Justus Claproth in 1774. Today this method is called deinking. It was not until the introduction of wood pulp in 1843 that paper production was not dependent on recycled materials from ragpickers.
Fiber crops
Bast fibers (stem)
Bamboo, when derived from a mechanical process.
Dogbane, used by Native Americans
Esparto, a fiber from a grass
Flax, from which linen is derived
Hemp, a soft, strong fiber, edible seeds
Hoopvine, also used for barrel hoops and baskets, edible leaves, medicine
Jute, widely used, it is the cheapest fiber after cotton
Kenaf, the interior of the plant stem is used for its fiber. Edible leaves.
Lotus, used to produce lotus silk
Nettles used to make thread and twine, clothing made from it is both durable yet soft
Papyrus, a pith fiber, akin to a bast fiber
Ramie, a member of the nettle family.
Spanish broom, a legume, its fiber has similar characteristics to linen.
Tilia, known as Linden or Lime in Europe and Basswood in North America. Fiber comes from inner bark.
Leaf fibers
Abacá, a banana, producing "manila" rope from leaves
Piña, from pineapple leaves
Sisal, an agave
Bowstring hemp, a common house plant, also Sansevieria roxburghiana, Sansevieria hyacinthoides
Henequen, an agave. A useful fiber, but not as high quality as sisal
Phormium, "New Zealand Flax"
Yucca, an agave relative
Seed fibers and fruit fibers
Coir, the fiber from the coconut husk
Cotton
Kapok
Milkweed, grown for the filament-like pappus in its seed pods
Luffa, a gourd which when mature produces a sponge-like mass of xylem, used to make loofa sponge.
Walissima is a natural plant fiber obtained from Sida rhombifolia of Malvaceae family. It is produced mainly in Philippines islands.
| Technology | Basics_2 | null |
519903 | https://en.wikipedia.org/wiki/Spaniel | Spaniel | A spaniel is a type of gun dog. Spaniels were especially bred to flush game out of denser brush. By the late 17th century, spaniels had been specialized into water and land breeds. The extinct English Water Spaniel was used to retrieve water fowl shot down with arrows. Land spaniels were setting spaniels—those that crept forward and pointed their game, allowing hunters to ensnare them with nets, and springing spaniels — those that sprang pheasants and partridges (for hunting with falcon) and also rabbits and smaller mammals such as rats and mice (for hunting with greyhounds). During the 17th century, the role of the spaniel dramatically changed as Englishmen began hunting with flintlocks for wing shooting. Charles Goodall and Julia Gasow (1984) write that spaniels were "transformed from untrained, wild beaters, to smooth, polished gun dogs."
The word "spaniel" would seem to be derived from the medieval French espaigneul"Spanish"to modern French, espagnol.
Definition and description
The Oxford English Dictionary defines Spaniel as "a breed of dog with a long silky coat and drooping ears".
Not much has changed about spaniels in general over the years, as can be seen in this 1921 entry in Collier's New Encyclopedia:
Their distinguishing characteristics are a rather broad muzzle, remarkably long and full ears, hair plentiful and beautifully waved, particularly that of the ears, tail, and hinder parts of the thighs and legs. The prevailing color is liver and white, sometimes red and white or black and white, and sometimes deep brown, or black on the face and breast, with a tan spot over each eye. The English spaniel is a superior and very pure breed. The King Charles is a small variety of the spaniel used as a lapdog. The water spaniels, large and small, differ from the common spaniel only in the roughness of their coats, and in uniting the aquatic propensities of the Newfoundland dog with the fine hunting qualities of their own race. Spaniels possess a great share of intelligence, affection, and obedience, which qualities, combined with much beauty, make them highly prized as companions.
History
The origin of the word spaniel is described by the Oxford English Dictionary as coming from the Old French word espaigneul which meant "Spanish (dog)"; this in turn originated from the Latin Hispaniolus which simply means "Spanish".
In Edward, 2nd Duke of York's work The Master of Game, which was mostly a 15th-century translation of an earlier work by Gaston III of Foix-Béarn entitled Livre de chasse, spaniels are described as being from Spain as much as all Greyhounds are from England or Scotland. Sixteenth-century English physician John Caius wrote that the spaniels of the time were mostly white, marked with spots that are commonly red. He described a new variety to have come out of France, which were speckled all over with white and black, "which mingled colours incline to a marble blewe".
Celtic origin theory
In the appendices added to the 1909 re-print of Caius' work, the editors suggested that the type of dogs may have been brought into the British Isles as early as 900 BC by a branch of the Celts moving from Spain into Cornwall and on into Wales, England and Ireland. Theories on the origin of the Welsh Springer Spaniel support this theory, as it is believed that the breed specifically is a direct descendant of the "Agassian hunting dog" described in the hunting poem Cynegetica attributed to Oppian of Apamea, which belonged to the Celtic tribes of Roman Britain:
There is a strong breed of hunting dog, small in size but no less worthy of great praise. These the wild tribes of Britons with their tattooed backs rear and call by the name of Agassian. Their size is like that of worthless and greedy domestic table dogs; squat, emaciated, shaggy, dull of eye, but endowed with feet armed with powerful claws and a mouth sharp with close-set venomous tearing teeth. It is by virtue of its nose, however, that the Agassian is most exalted, and for tracking it is the best there is; for it is very adept at discovering the tracks of things that walk upon the ground, and skilled too at marking the airborne scent.
Roman origin theory
Another theory of the origin of the spaniel is that the ancient Romans imported spaniels into Britannia by way of the trade routes to the Far East. Colonel David Hancock adds a belief that the sporting type of spaniel originated in China from the short-faced ancestors of dogs such as the Pekingese, Pug and Shih Tzu. The theory goes that these ancestors were introduced into Southern Europe and evolved into the small sporting spaniels of the period around AD 1300–1600. The issue of how a short-muzzled dog could evolve into a longer-muzzled dog is addressed by pointing to the evolution of the King Charles Spaniel into the Cavalier King Charles Spaniel in less than a century.
Hunting
In assisting hunters, it is desirable that spaniels work within gun range, are steady to shot, and are able to mark the fall and retrieve shot game to hand with a soft mouth. A good nose is highly valued, as it is in most gun dog breeds. They are versatile hunters traditionally being used for upland game birds, but are equally adept at hunting rabbits, waterfowl, rats, and mice. Whether hunting in open fields, woodlands, farm lands—in briars, along fencerows or marshlands, a spaniel can get the job done.
On the basis of function and hunting style, the Fédération Cynologique Internationale (FCI) draws a distinction between Continental and Anglo-American spaniels. The FCI places Continental dogs of the spaniel type in the pointing group (Group 7, sect. 1.2) because they function more like setters which "freeze" and point to game. Breeds in this group include the Blue Picardy Spaniel, the French Spaniel, the Brittany, the Pont-Audemer Spaniel, and the Small Münsterländer. The FCI classifies most other dogs of the spaniel type as flushing or water dogs (Group 8, sections 2 and 3).
Breeds
Contemporary
Extinct
Misnamed
The following breeds are not true spaniels, but are named as such due to their resemblance to the spaniels.
| Biology and health sciences | Dogs | Animals |
519924 | https://en.wikipedia.org/wiki/Marattiaceae | Marattiaceae | Marattiaceae is the only family of extant (living) ferns in the order Marattiales. In the Pteridophyte Phylogeny Group classification of 2016 (PPG I), Marattiales is the only order in the subclass Marattiidae. The family has six genera and about 110 species. Many are different in appearance from other ferns, having large fronds and fleshy rootstocks.
Description
The Marattiaceae diverged from other ferns very early in their evolutionary history and are quite different from many plants familiar to people in temperate zones. Many of them have massive, fleshy rootstocks and the largest known fronds of any fern. The Marattiaceae is one of two groups of ferns traditionally known as eusporangiate ferns, meaning that the sporangium is formed from a group of cells as opposed to a leptosporangium in which there is a single initial cell. At least two genera, Angiopteris and Marattia, have been reported to undergo monoplastidic meiosis rather than polyplastidic meiosis, and are the only known examples within euphyllophytes to do so.
The large fronds characteristic of the group are most readily found in the genus Angiopteris, native to Australasia, Madagascar and Oceania. These fronds may be up to 9 meters long in the species Angiopteris teysmanniana of Java. In the Hawaiian Islands, Costa Rica, and Jamaica, the species Angiopteris evecta is naturalized, having escaped from botanical gardens, and is considered an invasive species. Angiopteris is unique among ferns in having explosively dispersed spores, which may contribute to its ability to spread.
Marattia in the strict sense is found in the neotropics and Hawaii with six recognized species. The genus Eupodium is also neotropical, with three species, and was originally described for the distinctive stalked synangia of some species.
Ptisana is a paleotropical genus, formerly thought to be part of Marattia. These plants are 2-4 times pinnate, with fronds often comparable in size to those found in Angiopteris. Terminal segments usually have a prominent suture where they attach. The sporangia lack the labiate apertures of Marattia and Eupodium, and synangia are deeply cut. The name of the genus derives from the resemblance of the synangia to pearl barley. The king fern, Ptisana salicina, from New Zealand and the South Pacific and known in Māori as "para" now has been placed in this genus. Sometimes called the potato fern, this is a large fern with an edible fleshy rhizome that is used as a food source by some indigenous peoples.
The East-Asian genus Christensenia is named in honor of the Danish pteridologist Carl Christensen is an uncommon fern with distinctive fronds resembling a horse chestnut leaf, hence the species Christensenia aesculifolia, meaning horse-chestnut-leaved Christensenia. Despite the relatively diminutive size of plants in this genus, the stomata of Christensenia are the largest known in the plant kingdom.
The genus Danaea is endemic to the Neotropics. They have bipinnate leaves with opposite pinnae, which are dimorphic, the fertile leaves much contracted, and covered below with sunken, linear synangia dehiscing via pores.
Taxonomy
in the Pteridophyte Phylogeny Group classification of 2016 (PPG I), Marattiaceae is the only family in the order Marattiales, which in turn is the only order in the subclass Marattiidae. Marattiidae is one of four subclasses of class Polypodiopsida (ferns), to which it is related as shown in this cladogram, being a sister group to Polypodiidae.
History of classification
In the molecular phylogenetic classification of Smith et al. in 2006, the Marattiales formed the single member of the class Marattiopsida. Four genera were recognized. The class was lowered in rank to the subclass Marattiidae in the 2009 classification of Mark W. Chase and James L. Reveal, and subsequent systems such as Christenhusz et al. (2011). The Pteridophyte Phylogeny Group (2016) classification retains this rank. In that system, Marattiidae is monotypic and has one order, Marattiales, one family, Marattiaceae, six genera, and an estimated 111 species.
There have long been four traditional extant genera (Angiopteris, Christensenia, Danaea and Marattia), but phylogenetic analysis has determined the genus Marattia to be paraphyletic, and the genus has been split into three genera, Marattia in the strict sense, Eupodium, and Ptisana. Christenhusz and Chase placed Danaea in subfamily Danaeoideae and the remaining genera in subfamily Marattioideae, but this subfamilial classification was not taken up by PPG I.
This fern group has a long fossil history with many extinct taxa (Psaronius, Asterotheca, Scolecopteris, Eoangiopteris, Qasimia, Marantoidea, Danaeites, Marattiopsis, Ptychocarpus, etc.).
Genera
Six genera are accepted in the PPG I classification:
Angiopteris Hoffm.
Christensenia Maxon
Danaea Sm.
Eupodium J.Sm.
Marattia Sw.
Ptisana Murdock
Several other genera have been named in the Marattiaceae, namely: Archangiopteris, Clementea, Macroglossum, Protangiopteris, Protomarattia and Psilodochea. These are currently treated as synonyms of Angiopteris.
Evolutionary history
Marattiaceae are considered one of the most primitive living lineages of ferns. The earliest members of the family appeared during the Carboniferous, over 300 million years ago. The group has an extensive fossil record extending from the Carboniferous into the Jurassic, but post-Jurassic records are scarce.
| Biology and health sciences | Ferns | Plants |
520222 | https://en.wikipedia.org/wiki/Antarctic%20plate | Antarctic plate | The Antarctic plate is a tectonic plate containing the continent of Antarctica, the Kerguelen Plateau, and some remote islands in the Southern Ocean and other surrounding oceans. After breakup from Gondwana (the southern part of the supercontinent Pangea), the Antarctic plate began moving the continent of Antarctica south to its present isolated location, causing the continent to develop a much colder climate. The Antarctic plate is bounded almost entirely by extensional mid-ocean ridge systems. The adjoining plates are the Nazca plate, the South American plate, the African plate, the Somali plate, the Indo-Australian plate, the Pacific plate, and, across a transform boundary, the Scotia and South Sandwich plates.
The Antarctic plate has an area of about . It is Earth's fifth-largest tectonic plate.
The Antarctic plate's movement is estimated to be at least per year towards the Atlantic Ocean.
Subduction beneath South America
The Antarctic plate started to subduct beneath South America 14 million years ago in the Miocene epoch. At first it subducted only in the southernmost tip of Patagonia, meaning that the Chile triple junction lay near the Strait of Magellan. As the southern part of the Nazca plate and the Chile Rise became consumed by subduction the more northerly regions of the Antarctic plate began to subduct beneath Patagonia so that the Chile triple junction lies at present in front of Taitao Peninsula at 46°15' S.
The subduction of the Antarctic plate beneath South America is held to have uplifted Patagonia as it reduced the previously vigorous down-dragging flow in the Earth's mantle caused by the subduction of the Nazca plate beneath Patagonia. The dynamic topography caused by this uplift raised Quaternary-aged marine terraces and beaches across the Atlantic coast of Patagonia.
Land
Amsterdam Island (France)
Antarctica
East Antarctica
Transantarctic Mountains
West Antarctica
Antarctic Peninsula
Crozet Islands (France)
Heard Island and McDonald Islands (Australia)
Heard Island
McDonald Islands
Kerguelen Islands (France)
Peter I Island
Prince Edward Islands (South Africa)
Saint Paul Island (France)
South Orkney Islands
South Shetland Islands
| Physical sciences | Tectonic plates | Earth science |
520289 | https://en.wikipedia.org/wiki/Hearing%20aid | Hearing aid | A hearing aid is a device designed to improve hearing by making sound audible to a person with hearing loss. Hearing aids are classified as medical devices in most countries, and regulated by the respective regulations. Small audio amplifiers such as personal sound amplification products (PSAPs) or other plain sound reinforcing systems cannot be sold as "hearing aids".
Early devices, such as ear trumpets or ear horns, were passive amplification cones designed to gather sound energy and direct it into the ear canal.
Modern devices are computerised electroacoustic systems that transform environmental sound to make it audible, according to audiometrical and cognitive rules. Modern devices also utilize sophisticated digital signal processing, aiming to improve speech intelligibility and comfort for the user. Such signal processing includes feedback management, wide dynamic range compression, directionality, frequency lowering, and noise reduction.
Modern hearing aids require configuration to match the hearing loss, physical features, and lifestyle of the wearer. The hearing aid is fitted to the most recent audiogram and is programmed by frequency. This process called "fitting" can be performed by the user in simple cases, by a Doctor of Audiology, also called an audiologist (AuD), or by a Hearing Instrument Specialist (HIS) or audioprosthologist. The amount of benefit a hearing aid delivers depends in large part on the quality of its fitting. Almost all hearing aids in use in the US are digital hearing aids, as analog aids are phased out. Devices similar to hearing aids include the osseointegrated auditory prosthesis (formerly called the bone-anchored hearing aid) and cochlear implant.
Uses
Hearing aids are used for a variety of pathologies including sensorineural hearing loss, conductive hearing loss, and single-sided deafness. Hearing aid candidacy was traditionally determined by a Doctor of Audiology, or a certified hearing specialist, who will also fit the device based on the nature and degree of the hearing loss being treated. The amount of benefit experienced by the user of the hearing aid is multi-factorial, depending on the type, severity, and etiology of the hearing loss, the technology and fitting of the device, and on the motivation, personality, lifestyle, and overall health of the user. Over-the-counter hearing aids, which address mild to moderate hearing loss, are designed to be adjusted by the user.
Hearing aids are incapable of truly correcting a hearing loss; they are an aid to make sounds more audible. The most common form of hearing loss for which hearing aids are sought is sensorineural, resulting from damage to the hair cells and synapses of the cochlea and auditory nerve. Sensorineural hearing loss reduces the sensitivity to sound, which a hearing aid can partially accommodate by making sound louder. Other decrements in auditory perception caused by sensorineural hearing loss, such as abnormal spectral and temporal processing, and which may negatively affect speech perception, are more difficult to compensate for using digital signal processing and in some cases may be exacerbated by the use of amplification. Conductive hearing losses, which do not involve damage to the cochlea, tend to be better treated by hearing aids; the hearing aid is able to sufficiently amplify sound to account for the attenuation caused by the conductive component. Once the sound is able to reach the cochlea at normal or near-normal levels, the cochlea and auditory nerve are able to transmit signals to the brain normally.
Common issues with hearing aid fitting and use are the occlusion effect, loudness recruitment, and understanding speech in noise. Once a common problem, feedback is generally now well-controlled through the use of feedback management algorithms.
Candidacy and acquisition
There are several ways of evaluating how well a hearing aid compensates for hearing loss. One approach is audiometry which measures a subject's hearing levels in laboratory conditions. The threshold of audibility for various sounds and intensities is measured in a variety of conditions. Although audiometric tests may attempt to mimic real-world conditions, the patient's own every day experiences may differ. An alternative approach is self-report assessment, where the patient reports their experience with the hearing aid.
Hearing aid outcome can be represented by three dimensions:
hearing aid usage
aided speech recognition
benefit/satisfaction
The most reliable method for assessing the correct adjustment of a hearing aid is through real ear measurement. Real ear measurements (or probe microphone measurements) are an assessment of the characteristics of hearing aid amplification near the ear drum using a silicone probe tube microphone.
Current research is also pointing towards hearing aids and proper amplification as a treatment for tinnitus, a medical condition which manifests itself as a ringing or buzzing in the ears.
Types
There are many types of hearing aids (also known as hearing instruments), which vary in size, power and circuitry.
Among the different sizes and models are:
Body-worn
Body worn aids were the first portable electronic hearing aids, and were invented by Harvey Fletcher while working at Bell Laboratories. Body aids consist of a case and an earmold, attached by a wire. The case contains the electronic amplifier components, controls and battery, while the earmold typically contains a miniature loudspeaker. The case is typically about the size of a pack of playing cards and is carried in a pocket or on a belt.
Without the size constraints of smaller hearing devices, body worn aid designs can provide large amplification and long battery life at a lower cost. Body aids are still used in emerging markets because of their relatively low cost.
Behind the ear
Behind the ear hearing aids are one of two major classes of hearing aids – behind the ear (BTE) and in the ear (ITE). These two classes are distinguished by where the hearing aid is worn. BTE hearing aids consist of a case which hangs behind the pinna. The case is attached to an earmold or dome tip by a traditional tube, slim tube, or wire. The tube or wire courses from the superior-ventral portion of the pinna to the concha, where the ear mold or dome tip inserts into the external auditory canal. The case contains the electronics, controls, battery, and microphone(s).The loudspeaker, or receiver, may be housed in the case (traditional BTE) or in the earmold or dome tip (receiver-in-the-canal, or RIC). The RIC style of BTE hearing aid is often smaller than a traditional BTE and more commonly used in more active populations.
BTEs are generally capable of providing more output and may therefore be indicated for more severe degrees of hearing loss. However, BTEs are very versatile and can be used for nearly any kind of hearing loss. BTEs come in a variety of sizes, ranging from a small, "mini BTE", to larger, ultra-power devices. Size typically depends on the output level needed, the location of the receiver, and the presence or absence of a telecoil. BTEs are durable, easy to repair, and often have controls and battery doors that are easier to manipulate. BTEs are also easily connected to assistive listening devices, such as FM systems and induction loops. BTEs are commonly worn by children who need a durable type of hearing aid.
In the ear
In the ear aids (ITE) devices fit in the outer ear bowl (called the concha). Being larger, these are easier to insert and can hold extra features. They are sometimes visible when standing face to face with someone. ITE hearing aids are custom made to fit each individual's ear. They can be used in mild to some severe hearing losses. Feedback, a squealing/whistling caused by sound (particularly high frequency sound) leaking and being amplified again, may be a problem for severe hearing losses. Some modern circuits are able to provide feedback regulation or cancellation to assist with this.
Venting may also cause feedback. A vent is a tube primarily placed to offer pressure equalization. However, different vent styles and sizes can be used to influence and prevent feedback.
Traditionally, ITEs have not been recommended for young children because their fit could not be as easily modified as the earmold for a BTE, and thus the aid had to be replaced frequently as the child grew. However, there are new ITEs made from a silicone type material that mitigates the need for costly replacements.
ITE hearing aids can be connected wirelessly to FM systems, for instance with a body-worn FM receiver with induction neck-loop which transmits the audio signal from the FM transmitter inductively to the telecoil inside the hearing instrument.
Mini in canal (MIC) or completely in canal (CIC) aids are generally not visible unless the viewer looks directly into the wearer's ear. These aids are intended for mild to moderately severe losses. CICs are usually not recommended for people with good low-frequency hearing, as the occlusion effect is much more noticeable. Completely-in-the-canal hearing aids fit tightly deep in the ear. It is barely visible. Being small, it will not have a directional microphone, and its small batteries will have a short life, and the batteries and controls may be difficult to manage. Its position in the ear prevents wind noise and makes it easier to use phones without feedback. In-the-canal hearing aids are placed deep in the ear canal. They are barely visible. Larger versions of these can have directional microphones. Being in the canal, they are less likely to cause a plugged feeling. These models are easier to manipulate than the smaller completely in-the-canal models but still have the drawbacks of being rather small.
In-the-ear hearing aids are typically more expensive than behind-the-ear counterparts of equal functionality, because they are custom fitted to the patient's ear.
In fitting, the audiologist takes a physical impression (mold) of the ear.
The mold is scanned by a specialized CAD system, resulting in a 3D model of the outer ear.
During modeling, the venting tube is inserted.
The digitally modeled shell is printed using a rapid prototyping technique such as stereolithography.
Finally, the aid is assembled and shipped to the audiologist after a quality check.
Invisible-in-canal hearing aids
Invisible-in-canal hearing aids (IIC) style of hearing aids fits inside the ear canal completely, leaving little to no trace of an installed hearing aid visible. This is because it fits deeper in the canal than other types, so that it is out of view even when looking directly into the ear bowl (concha). A comfortable fit is achieved because the shell of the aid is custom-made to the individual ear canal after taking a mold.
Invisible hearing aid types use venting and their deep placement in the ear canal to give a more natural experience of hearing. Unlike other hearing aid types, with the IIC aid the majority of the ear is not blocked (occluded) by a large plastic shell. This means that sound can be collected more naturally by the shape of the ear, and can travel down into the ear canal as it would with unassisted hearing. Depending on their size, some models allow the wearer to use a mobile phone as a remote control to alter memory and volume settings, instead of taking the IIC out to do this. IIC types are most suitable for users up to middle age, but are not suitable for elderly people with unsteady hands.
Extended wear hearing aids
Extended wear hearing aids are hearing devices that are non-surgically placed in the ear canal by a hearing professional. The extended wear hearing aid represents the first "invisible" hearing device. These devices are worn for 1–3 months at a time without removal. They are made of soft material designed to contour to each user and can be used by people with mild to moderately severe hearing loss. Their close proximity to the ear drum results in improved sound directionality and localization, reduced feedback, and improved high frequency gain. While traditional BTE or ITC hearing aids require daily insertion and removal, extended wear hearing aids are worn continuously and then replaced with a new device. Users can change volume and settings without the aid of a hearing professional. The devices are very useful for active individuals because their design protects against moisture and earwax and can be worn while exercising, showering, etc. Because the device's placement within the ear canal makes them invisible to observers, extended wear hearing aids are popular with those who are self-conscious about the aesthetics of BTE or ITC hearing aid models. As with other hearing devices, compatibility is based on an individual's hearing loss, ear size and shape, medical conditions, and lifestyle. The disadvantages include regular removal and reinsertion of the device when the battery dies, inability to go underwater, earplugs when showering, and for some discomfort with the fit since it is inserted deeply in the ear canal, the only part of the body where skin rests directly on top of bone.
CROS hearing aid
A CROS hearing aid is a hearing aid that transmits auditory information from one side of the head to the other side of the head. Candidates include people who have poor word understanding on one side, no hearing on one side, or who are not benefiting from a hearing aid on one side. CROS hearing aids can appear very similar to behind the ear hearing aids. The CROS system can assist the patient in sound localization and understanding auditory information on their poor side. While CROS hearing aids can be quite effective, the long-term solution for those with hearing issues on one side is to use a BiCROS system. This creates more of a balance for wearers.
Bone-anchored
A bone anchored hearing aid (BAHA) is a surgically implanted auditory prosthetic based on bone conduction. It is an option for patients without external ear canals, when conventional hearing aids with a mold in the ear cannot be used. The BAHA uses the skull as a pathway for sound to travel to the inner ear. For people with conductive hearing loss, the BAHA bypasses the external auditory canal and middle ear, stimulating the functioning cochlea. For people with unilateral hearing loss, the BAHA uses the skull to conduct the sound from the deaf side to the side with the functioning cochlea.
Individuals under the age of two (five in the USA) typically wear the BAHA device on a Softband. This can be worn from the age of one month as babies tend to tolerate this arrangement very well. When the child's skull bone is sufficiently thick, a titanium "post" can be surgically embedded into the skull with a small abutment exposed outside the skin. The BAHA sound processor sits on this abutment and transmits sound vibrations to the external abutment of the titanium implant. The implant vibrates the skull and inner ear, which stimulate the nerve fibers of the inner ear, allowing hearing.
The surgical procedure is simple both for the surgeon, involving very few risks for the experienced ear surgeon. For the patient, minimal discomfort and pain is reported. Patients may experience numbness of the area around the implant as small superficial nerves in the skin are sectioned during the procedure. This often disappears after some time. There is no risk of further hearing loss due to the surgery. One important feature of the BAHA is that, if a patient for whatever reason does not want to continue with the arrangement, it takes the surgeon less than a minute to remove it. The BAHA does not restrict the wearer from any activities such as outdoor life, sporting activities etc.
A BAHA can be connected to an FM system by attaching a miniaturized FM receiver to it.
Two main brands manufacture BAHAs today – the original inventors Cochlear, and the hearing aid company Oticon.
Eyeglass aids
During the late 1950s through 1970s, before in-the-ear aids became common (and in an era when thick-rimmed eyeglasses were popular), people who wore both glasses and hearing aids frequently chose a type of hearing aid that was built into the
temple pieces of the spectacles. However, the combination of glasses and hearing aids was inflexible: the range of frame styles was limited, and the user had to wear both hearing aids and glasses at once or wear neither. Today, people who use both glasses and hearing aids can use in-the-ear types, or rest a BTE neatly alongside the arm of the glasses. There are still some specialized situations where hearing aids built into the frame of eyeglasses can be useful, such as when a person has hearing loss mainly in one ear: sound from a microphone on the "bad" side can be sent through the frame to the side with better hearing.
This can also be achieved by using CROS or bi-CROS style hearing aids, which are now wireless in sending sound to the better side.
Spectacle hearing aids
These are generally worn by people with a hearing loss who either prefer a more cosmetic appeal of their hearing aids by being attached to their glasses or where sound cannot be passed in the normal way, via a hearing aids, perhaps due to a blockage in the ear canal. pathway or if the client experiences continual infections in the ear.
Spectacle aids come in two forms, bone conduction spectacles and air conduction spectacles.
Bone conduction spectacles
Sounds are transmitted via a receiver attached from the arm of the spectacles which are fitted firmly behind the boney portion of the skull at the back of the ear, (mastoid process) by means of pressure, applied on the arm of the spectacles. The sound is passed from the receiver on the arm of the spectacles to the inner ear (cochlea), via the bony portion. The process of transmitting the sound through the bone requires a great amount of power. Bone conduction aids generally have a poorer high pitch response and are therefore best used for conductive hearing losses or where it is impractical to fit standard hearing aids.
Air conduction spectacles
Unlike the bone conduction spectacles the sound is transmitted via hearing aids which are attached to the arm or arms of the spectacles. When removing your glasses for cleaning, the hearing aids are detached at the same time. Whilst there are genuine instances where spectacle aids are a preferred choice, they may not always be the most practical option.
Directional spectacles
These 'hearing glasses' incorporate a directional microphone capability: four microphones on each side of the frame effectively work as two directional microphones, which are able to discern between sound coming from the front and sound coming from the sides or back of the user. This improves the signal-to-noise ratio by allowing for amplification of the sound coming from the front, the direction in which the user is looking, and active noise control for sounds coming from the sides or behind.
Only very recently has the technology required become small enough to be fitted in the frame of the glasses. As a recent addition to the market, this new hearing aid is currently available only in the Netherlands and Belgium.
Stethoscope
These hearing aids are designed for medical practitioners with hearing loss who use stethoscopes. The hearing aid is built into the speaker of the stethoscope, which amplifies the sound.
Hearing aid applications
Hearing aid applications (HAA) are software which, when installed on mobile computational platforms, transforms them into hearing aids.
The principle of HAA operation corresponds to the basic principles of operation of traditional hearing aids: the microphone receives an acoustic signal and converts it into a digital form. Sound amplification is achieved by the means of a mobile computational platform, in accordance with the degree and type of the user's hearing loss. The processed audio signal is transformed into an audio signal and output to the user into the headphones/headset. Signal processing is implemented in real time.
Constructional features of mobile computational platforms imply preferred use of stereo headsets with two speakers, which allows carrying out binaural hearing correction for the left and right ear separately. HAAs can work with both wired and wireless headsets and headphones.
As a rule, HAAs have two operation modes: setup mode and hearing aid mode. Setup mode involves the user passing an in situ-audiometry procedure, which determines the user's hearing characteristics. Hearing aid mode is a hearing correction system that corrects the user's hearing in accordance with the user's hearing thresholds. HAAs also incorporate background noise suppression and acoustic feedback suppression.
The user can independently choose a formula to enhance the sound, as well as adjust the level of the desired amplification to their wishes.
HAAs have several advantages (compared to traditional hearing aids):
HAAs do not cause any psychological inconvenience;
it is possible to achieve the highest sound pressure level and get high sound quality (due to large speakers and a long battery life);
it is possible to use more complex audio signal processing algorithms and a higher sampling rate (because of capacious battery);
the possibility to implement more convenient application control functions for people with poor motor skills;
resistance to ingress of earwax and moisture;
software flexibility;
the large distance between the microphone and the speaker prevents the occurrence of acoustic feedback;
the set up of HAAs in simple cases does not require special equipment and qualifications;
the user does not need to purchase and carry any separate device;
various types of headphones and headsets can be used.
HAAs also have some disadvantages (compared to traditional hearing aids):
because the microphone is not located in the ear, it does not use the functional advantages of the auricle and the natural acoustics of the outer ear.
they are more noticeable and less comfortable to wear.
Technology
The first electrical hearing aid used the carbon microphone of the telephone and was introduced in 1896. The vacuum tube made electronic amplification possible, but early versions of amplified hearing aids were too heavy to carry around. Miniaturization of vacuum tubes lead to portable models, and after World War II, wearable models using miniature tubes. The transistor invented in 1948 was well suited to the hearing aid application due to low power and small size; hearing aids were an early adopter of transistors. The development of integrated circuits allowed further improvement of the capabilities of wearable aids, including implementation of digital signal processing techniques and programmability for the individual user's needs.
Compatibility with telephones
A hearing aid and a telephone are "compatible" when they can connect to each other in a way that produces clear, easily understood sound. The term "compatibility" is applied to all three types of telephones (wired, cordless, and mobile). There are two ways telephones and hearing aids can connect with each other:
Acoustically: the sound from the phone's speaker is picked up by the hearing aid's microphone.
Electromagnetically: the signal inside the phone's speaker is picked up by the hearing aid's "telecoil" or "T-coil", a special loop of wire inside the hearing aid.
Note that telecoil coupling has nothing to do with the radio signal in a cellular or cordless phone: the audio signal picked up by the telecoil is the weak electromagnetic field that is generated by the voice coil in the phone's speaker as it pushes the speaker cone back and forth.
The electromagnetic (telecoil) mode is usually more effective than the acoustic method. This is mainly because the microphone is often automatically switched off when the hearing aid is operating in telecoil mode, so background noise is not amplified. Since there is an electronic connection to the phone, the sound is clearer and distortion is less likely. But in order for this to work, the phone has to be hearing-aid compatible. More technically, the phone's speaker has to have a voice coil that generates a relatively strong electromagnetic field. Speakers with strong voice coils are more expensive and require more energy than the tiny ones used in many modern telephones; phones with the small low-power speakers cannot couple electromagnetically with the telecoil in the hearing aid, so the hearing aid must then switch to acoustic mode. Also, many mobile phones emit high levels of electromagnetic noise that creates audible static in the hearing aid when the telecoil is used. A workaround that resolves this issue on many mobile phones is to plug a wired (not Bluetooth) headset into the mobile phone; with the headset placed near the hearing aid the phone can be held far enough away to attenuate the static. Another method is to use a "neckloop" (which is like a portable, around-the-neck induction loop), and plug the neckloop directly into the standard audio jack (headphones jack) of a smartphone (or laptop, or stereo, etc.). Then, with the hearing aids' telecoil turned on (usually a button to press), the sound will travel directly from the phone, through the neckloop and into the hearing aids' telecoils.
On 21 March 2007, the Telecommunications Industry Association issued the TIA-1083 standard, which gives manufacturers of cordless telephones the ability to test their products for compatibility with most hearing aids that have a T-Coil magnetic coupling mode. With this testing, digital cordless phone manufacturers will be able to inform consumers about which products will work with their hearing aids.
The American National Standards Institute (ANSI) has a ratings scale for compatibility between hearing aids and phones:
When operating in acoustic (Microphone) mode, the ratings are from M1 (worst) to M4 (best).
When operating in electromagnetic (Telecoil) mode, the ratings are from T1 (worst) to T4 (best).
The best possible rating is M4/T4 meaning that the phone works well in both modes. Devices rated below M3 are unsatisfactory for people with hearing aids.
Computer programs that allow the creation of a hearing aid using a PC, tablet or smartphone are currently gaining in popularity. Modern mobile devices have all the necessary components to implement this: hardware (an ordinary microphone and headphones may be used) and a high-performance microprocessor that carries digital sound processing according to a given algorithm.
Application configuration is carried out by the user themselves in accordance with the individual features of their hearing ability. The computational power of modern mobile devices is sufficient to produce the best sound quality. This, coupled with software application settings (for example, profile selection according to a sound environment) provides for high comfort and convenience of use.
In comparison with the digital hearing aid, mobile applications have the following advantages:
acoustic gain is up to 30 dB (with a standard headset);
complete invisibility (smartphone is not associated with a hearing aid);
ease of use (no need to use additional devices, batteries and so on.);
Fast switching between the external headset and phone microphone;
free distribution of applications.
High duration of the battery;
high sampling frequency (44.1 kHz) providing for excellent sound quality;
high wearing comfort;
low delay in audio processing (from 6,3 to 15,7 ms – depending on the mobile device model);
No loss of settings when switching from one gadget to another and back again;
No need to get used to it, when changing mobile devices;
user-friendly interface of software settings;
It should be clearly understood that "hearing aid" application for smartphone / tablet cannot be considered a complete substitution of a digital hearing aid, since the latter:
is a medical device (exposed to the relevant procedures of testing and certification);
is adjusted using audiometry procedures.
is designed for use by doctor's prescription;
Functionality of hearing aid applications may involve a hearing test (in situ audiometry) too. However, the results of the test are used only to adjust the device for comfortable working with the application. The procedure of hearing testing in any way cannot claim to replace an audiometry test carried out by a medical specialist, so cannot be a basis for diagnosis.
Apps such as Oticon ON for certain iOS (Apple) and Android devices can assist in locating a lost/misplaced hearing aid.
Wireless
Recent hearing aids include wireless hearing aids. One hearing aid can transmit to the other side so that pressing one aid's program button simultaneously changes the other aid, so that both aids change background settings simultaneously. FM listening systems are now emerging with wireless receivers integrated with the use of hearing aids. A separate wireless microphone can be given to a partner to wear in a restaurant, in the car, during leisure time, in the shopping mall, at lectures, or during religious services. The voice is transmitted wirelessly to the hearing aids eliminating the effects of distance and background noise. FM systems have shown to give the best speech understanding in noise of all available technologies.
FM systems can also be hooked up to a TV or a stereo.
2.4 gigahertz Bluetooth connectivity is the most recent innovation in wireless interfacing for hearing instruments to audio sources such as TV streamers or Bluetooth enabled mobile phones. Current hearing aids generally do not stream directly via Bluetooth but rather do so through a secondary streaming device (usually worn around the neck or in a pocket), this bluetooth enabled secondary device then streams wirelessly to the hearing aid but can only do so over a short distance. This technology can be applied to ready-to-wear devices (BTE, Mini BTE, RIE, etc.) or to custom made devices that fit directly into the ear.
In developed countries FM systems are considered a cornerstone in the treatment of hearing loss in children. More and more adults discover the benefits of wireless FM systems as well, especially since transmitters with different microphone settings and Bluetooth for wireless cell phone communication have become available.
Many theatres and lecture halls are now equipped with assistive listening systems that transmit the sound directly from the stage; audience members can borrow suitable receivers and hear the program without background noise. In some theatres and churches FM transmitters are available that work with the personal FM receivers of hearing instruments.
Directional microphone
Most older hearing aids have only an omnidirectional microphone. An omnidirectional microphone amplifies sounds equally from all directions. In contrast, a directional microphone amplifies sounds from one direction more than sounds from other directions. This means that sounds originating from the direction the system is steered toward are amplified more than sounds coming from other directions. If the desired speech arrives from the direction of steering and the noise is from a different direction, then compared to an omnidirectional microphone, a directional microphone provides a better signal-to-noise ratio. Improving the signal-to-noise ratio improves speech understanding in noise. Directional microphones have been found to be the second best method to improve the signal-to-noise ratio (the best method was an FM system, which locates the microphone near the mouth of the desired talker).
Many hearing aids have both an omnidirectional and a directional microphone mode. This is because the wearer may not need or desire the noise-reducing properties of the directional microphone in a given situation. Typically, the omnidirectional microphone mode is used in quiet listening situations (e.g. living room) whereas the directional microphone is used in noisy listening situations (e.g. restaurant). The microphone mode is typically selected manually by the wearer. Some hearing aids automatically switch the microphone mode.
Adaptive directional microphones automatically vary the direction of maximum amplification or rejection (to reduce an interfering directional sound source). The direction of amplification or rejection is varied by the hearing aid processor. The processor attempts to provide maximum amplification in the direction of the desired speech signal source or rejection in the direction of the interfering signal source. Unless the user manually temporarily switches to a "restaurant program, forward only mode" adaptive directional microphones frequently amplify the speech of other talkers in a cocktail party type environments, such as restaurants or coffee shops; this can also be helpful during business meetings. The presence of multiple speech signals makes it difficult for the processor to correctly select the desired speech signal. Another disadvantage is that some noises often contain characteristics similar to speech, making it difficult for the hearing aid processor to distinguish the speech from the noise. Despite the disadvantages, adaptive directional microphones can provide improved speech recognition in noise.
FM systems have been found to provide a better signal-to-noise ratio even at larger speaker-to-talker distances in simulated testing conditions.
Telecoil
Telecoils or T-coils (from "Telephone Coils") are small devices installed in hearing aids or cochlear implants. An audio induction loop generates an electromagnetic field that can be detected by T-coils, allowing audio sources to be directly connected to a hearing aid. The T-coil is intended to help the wearer filter out background noise. They can be used with telephones, FM systems (with neck loops), and induction loop systems (also called "hearing loops") that transmit sound to hearing aids from public address systems and TVs. In the UK and the Nordic countries, hearing loops are widely used in churches, shops, railway stations, and other public places. In the US, telecoils and hearing loops are gradually becoming more common. Audio induction loops, telecoils and hearing loops are gradually becoming more common also in Slovenia.
A T-coil consists of a metal core (or rod) around which ultra-fine wire is coiled. T-coils are also called induction coils because when the coil is placed in a magnetic field, an alternating electric current is induced in the wire (Ross, 2002b; Ross, 2004). The T-coil detects magnetic energy and transduces (converts) it to electrical energy. In the United States, the Telecommunications Industry Association's TIA-1083 standard, specifies how analog handsets can interact with telecoil devices, to ensure the optimal performance.
Although T-coils are effectively a wide-band receiver, interference is unusual in most hearing loop situations. Interference can manifest as a buzzing sound, which varies in volume depending on the distance the wearer is from the source. Sources are electromagnetic fields, such as CRT computer monitors, older fluorescent lighting, some dimmer switches, many household electrical appliances and airplanes.
The states of Florida and Arizona have passed legislation that requires hearing professionals to inform patients about the usefulness of telecoils.
Legislation affecting use
In the United States, the Hearing Aid Compatibility Act of 1988 requires that the Federal Communications Commission (FCC) ensure that all telephones manufactured or imported for use in the United States after August 1989, and all "essential" telephones, be hearing aid-compatible (through the use of a telecoil).
"Essential" phones are defined as "coin-operated telephones, telephones provided for emergency use, and other telephones frequently needed for use by persons using such hearing aids." These might include workplace telephones, telephones in confined settings (like hospitals and nursing homes), and telephones in hotel and motel rooms. Secure telephones, as well as telephones used with public mobile and private radio services, are exempt from the HAC Act. "Secure" phones are defined as "telephones that are approved by the U.S. Government for the transmission of classified or sensitive voice communications."
In 2003, the FCC adopted rules to make digital wireless telephones compatible with hearing aids and cochlear implants. Although analog wireless phones do not usually cause interference with hearing aids or cochlear implants, digital wireless phones often do because of electromagnetic energy emitted by the phone's antenna, backlight, or other components. The FCC has set a timetable for the development and sale of digital wireless telephones that are compatible with hearing aids. This effort promises to increase the number of digital wireless telephones that are hearing aid-compatible. Older generations of both cordless and mobile phones used analog technology.
Audio boot
An audio boot or audio shoe is an electronic device used with hearing aids; hearing aids often come with a special set of metal contacts for audio input. Typically the audio boot will fit around the end of the hearing aid (a behind-the-ear model, as in-the-ear do not afford any purchase for the connection) to link it with another device, like an FM system or a cellphone or even a digital audio player.
Direct audio input
Direct audio input (DAI) allows the hearing aid to be directly connected to an external audio source like a CD player or an assistive listening device (ALD). By its very nature, DAI is susceptible to far less electromagnetic interference, and yields a better quality audio signal as opposed to using a T-coil with standard headphones. An audio boot is a type of device that may be used to facilitate DAI.
Processing
Every electronic hearing aid has at minimum a microphone, a loudspeaker (commonly called a receiver), a battery, and electronic circuitry. The electronic circuitry varies among devices, even if they are the same style. The circuitry falls into three categories based on the type of audio processing (analog or digital) and the type of control circuitry (adjustable or programmable). Hearing aid devices generally do not contain processors strong enough to process complex signal algorithms for sound source localization.
Analog
Analog audio may have:
Adjustable control: The audio circuit is analog with electronic components that can be adjusted. The hearing professional determines the gain and other specifications required for the wearer, and then adjusts the analog components either with small controls on the hearing aid itself or by having a laboratory build the hearing aid to meet those specifications. After the adjustment the resulting audio does not change any further, other than overall loudness that the wearer adjusts with a volume control. This type of circuitry is generally the least flexible. The first practical electronic hearing aid with adjustable analog audio circuitry was based on US Patent 2,017,358, "Hearing Aid Apparatus and Amplifier" by Samual Gordon Taylor, filed in 1932.
Programmable control: The audio circuit is analog but with additional electronic control circuitry that can be programmed by an audiologist, often with more than one program. The electronic control circuitry can be fixed during manufacturing or in some cases, the hearing professional can use an external computer temporarily connected to the hearing aid to program the additional control circuitry. The wearer can change the program for different listening environments by pressing buttons either on the device itself or on a remote control or in some cases the additional control circuitry operates automatically. This type of circuitry is generally more flexible than simple adjustable controls. The first hearing aid with analog audio circuitry and automatic digital electronic control circuitry was based on US Patent 4,025,721, "Method of and means for adaptively filtering near-stationary noise from speech" by D Graupe, GD Causey, filed in 1975. This digital electronic control circuitry was used to identify and automatically reduce noise in individual frequency channels of the analog audio circuits and was known as the Zeta Noise Blocker.
Digital
Digital audio, programmable control: Both the audio circuit and the additional control circuits are fully digital. The hearing professional programs the hearing aid with an external computer temporarily connected to the device and can adjust all processing characteristics on an individual basis. Fully digital circuitry allows implementation of many additional features not possible with analog circuitry, can be used in all styles of hearing aids and is the most flexible; for example, digital hearing aids can be programmed to amplify certain frequencies more than others, and can provide better sound quality than analog hearing aids. Fully digital hearing aids can be programmed with multiple programs that can be invoked by the wearer, or that operate automatically and adaptively. These programs reduce acoustic feedback (whistling), reduce background noise, detect and automatically accommodate different listening environments (loud vs. soft, speech vs. music, quiet vs. noisy, etc.), control additional components such as multiple microphones to improve spatial hearing, transpose frequencies (shift high frequencies that a wearer may not hear to lower frequency regions where hearing may be better), and implement many other features. Fully digital circuitry also allows control over wireless transmission capability for both the audio and the control circuitry. Control signals in a hearing aid on one ear can be sent wirelessly to the control circuitry in the hearing aid on the opposite ear to ensure that the audio in both ears is either matched directly or that the audio contains intentional differences that mimic the differences in normal binaural hearing to preserve spatial hearing ability. Audio signals can be sent wirelessly to and from external devices through a separate module, often a small device worn like a pendant and commonly called a "streamer", that allows wireless connection to yet other external devices. This capability allows optimal use of mobile telephones, personal music players, remote microphones and other devices. With the addition of speech recognition and internet capability in the mobile phone, the wearer has optimal communication ability in many more situations than with hearing aids alone. This growing list includes voice activated dialing, voice activated software applications either on the phone or on the internet, receipt of audio signals from databases on the phone or on internet, or audio signals from television sets or from global positioning systems. The first practical, wearable, fully digital hearing aid was invented by Maynard Engebretson, Robert E Morley Jr. and Gerald R Popelka. Their work resulted in US Patent 4,548,082, "Hearing aids, signal supplying apparatus, systems for compensating hearing deficiencies, and methods" by A Maynard Engebretson, Robert E Morley Jr. and Gerald R Popelka, filed in 1984. This patent formed the basis of all subsequent fully digital hearing aids from all manufacturers, including those produced currently.
The signal processing is performed by the microprocessor in real time and taking into account the individual preferences of the user (for example, increasing bass for better speech perception in noisy environments, or selective amplification of high frequencies for people with reduced sensibility to this range). The microprocessor automatically analyzes the nature of the external background noise and adapts the signal processing to the specific conditions (as well as to its change, for example, when the user goes outside from the building).
In speech enhancement, for example using neural networks, finds application in hearing aids. Problems may arise if these methods filter out emergency sounds such as fire alarms and car horns.
Difference between digital and analog hearing aids
Analogue hearing aids make all the sounds picked up by the microphone louder. For example, speech and ambient noise will be made louder together. On the other hand, digital hearing aid (DHA) technology processes the sound using digital technology. Before transmitting the sound to the speaker, the DHA microprocessor processes the digital signal received by the microphone according to an algorithm. This allows certain-frequency sounds to be made louder according to the individual user's settings (personal audiogram), and the DHA can automatically adjust to various environments (noisy streets, quiet room, concert hall, etc.).
For users with varying degrees of hearing loss, it is difficult to perceive the entire frequency range of external sounds. DHAs with multi-channel digital processing allow a user to "compose" the output sound by fitting a whole spectrum of the input signal into it. This gives users with limited hearing abilities the opportunity to perceive the whole range of ambient sounds, despite the personal difficulties of perception of certain frequencies. Moreover, even in this "narrow" range the DHA microprocessor is able to emphasize desired sounds (e.g. speech), lowering unwanted loud, high, etc., sounds at the same time.
According to research DHAs have a number of significant advantages compared to analogue hearing aids:
"Self-learning" and adaptive adjustment. They can implement adaptive selection of amplification parameters and processing.
Effective acoustic feedback reduction. The acoustic whistling common to all hearing aids can be adaptively controlled.
Effective use of directional microphones. Directional microphones can be adaptively controlled.
Extended frequency range. A larger range of frequencies can be implemented with frequency shifting.
Flexibility in selective amplification. They can provide more flexibility in frequency specific amplification to match the individual hearing characteristics of the user.
Improved connection to other devices. Connection to other devices such as smartphones and televisions is possible.
Noise reduction. They can reduce the background noise level to increase user comfort in noisy environments.
Speech recognition. They can distinguish the speech signal from the overall spectrum of sounds, which facilitates speech perception.
These advantages of DHAs were confirmed by a number of studies relating to the comparative analysis of digital hearing aids of second and first generations and analog hearing aids.
Difference between digital hearing aids and hearing aid applications
Smartphones have all the necessary hardware to perform the functions of a digital hearing aid: microphone, AD converter, digital processor, DA converter, amplifier, and speakers. External microphone and speakers can also be connected as a special headset.
The operational principles of hearing aid applications correspond to general operational principles of digital hearing aids: the microphone perceives an acoustic signal and converts it to digital form. Sound amplification is achieved through hardware and software in accordance with the user's hearing characteristics. Then, the signal is converted to analog form and received in the headphones by the user. The signal is processed in real time.
Stereo headsets with two speakers can be used, which allows separate binaural hearing correction for the left and right ear.
Unlike digital hearing aids, the adjustment of hearing aid applications is an integral part of the application itself. Hearing aid applications are adjusted in accordance with the user's audiogram. The whole adjustment process is automated so that the user can perform audiometry on their own.
The hearing correction application has two modes: audiometry and correction. In the audiometry mode, hearing thresholds are measured. In the correction mode, the signal is processed with respect to the obtained thresholds.
Hearing aid applications also provide for different computational formulas for the calculation of sound amplification based on the audiometry data. These formulas are intended for maximum comfortable speech amplification and best sound intelligibility.
Hearing aid applications allow the user to save different user profiles for different acoustic environments. Thus, in contrast to the static settings of digital hearing aids, the user can quickly switch between the profiles depending on the acoustic environment.
One of the most important characteristics of the hearing aid is acoustic feedback. In hearing aid applications, there is a significant hardware delay, so hearing aid applications use a signal processing scheme with the minimum possible algorithmic delay to make it as short as possible.
Difference between PSAP and digital hearing aids
Personal sound amplification products (PSAP) are classified by the FDA as "personal sound amplification devices". These compact electronic devices are designed for people without hearing loss. Unlike hearing aids (which the FDA classifies as devices to compensate for hearing impairment), the use of PSAP does not require a medical prescription. Such devices are used by hunters, naturalists (for audio observation of animals or birds), ordinary people (for example, to increase the volume of the TV in a quiet room), etc. PSAP models differ significantly in price and functionality. Some devices simply amplify sound. Others contain directional microphones, equalizers to adjust the audio signal gain and filter noise. In modern days, some people refer to these devices as OTC hearing aids.
Evolution of hearing aid applications
There are audio players designed specifically for the hard-of-hearing. These applications amplify the volume of the reproduced audio signal in accordance with the user's hearing characteristics and act as a music volume amplifier and assistive hearing aid. The amplification algorithm works on the frequencies that the user hears worse, thus restoring natural hearing perception of the sound of music.
Just as in hearing aid applications, the player adjustment is based on the user's audiogram.
There are also applications that not only adapt the sound of music but also include some hearing aid functions. Such applications include a sound amplification mode in accordance with the user's hearing characteristics as well as a noise suppression mode and a mode allowing the user to hear ambient sound without pausing the music.
Also, some applications allow the hard-of-hearing to watch video and listen to the radio with comfort. The operational principles of these applications are similar to those of hearing aid applications: the audio signal is amplified on the frequencies that the user hears worse.
Hearing aid adaptation
A person using a hearing aid for the first time often cannot make use of all its advantages quickly. The structure and characteristics of hearing aids are thoroughly devised by specialists in order to make the adjustment period as simple and quick as possible. However, despite this, a beginning hearing aid user certainly needs time to get used to it.
The process of adjusting to hearing prostheses consists of the following steps:
Initial adjustment of the device
Fine adjustments
Adaptation to the new sound
Due to the plasticity of the central nervous system, inactive hearing centers in the brain's cortex switch over to processing auditory stimuli in another frequency and intensity. The brain starts to perceive sounds amplified by the hearing aid immediately after the initial adjustment; however, it may not process them correctly right away.
Feeling the hearing aid in the ear may seem unusual. It also takes time to adapt to a new way of hearing. The ear has to be gradually adjusted to the new sound. The sound may seem unnatural, metallic, too loud or too quiet. A whistling sound may also appear, which can be unpleasant.
Hearing aids do not provide immediate improvement. The adjustment period can last from several hours to several months.
Patients are offered an initial schedule to wear their hearing aid, ensuring gradual adaptation to it. If the patient wears the hearing aid continually from the beginning, the unfamiliar sound may cause a headache, and as a result, the user may refuse to wear a hearing aid despite the fact that it helps. Audiologists often run a quick preparation course for the patients. As a rule, users have inflated expectations of hearing aids. They expect that hearing aids will help them to hear in the same way as before hearing loss, but it is not the case. Training sessions help hearing aid users to get accustomed to the feeling of new sounds. Users are strongly recommended to regularly visit an audiologist, including for the purposes of additional hearing aid adjustment.
Hearing aid applications, in contrast to traditional hearing aids, allow the implementation of options such as a built-in adaptation course.
The functions of the course may include:
control of the amount of time spent on learning;
control over the sequence of exercises;
daily reminders to do the exercises.
The goal of the course is to help a user adapt to using a hearing aid application.
The adaptation course includes a certain number of stages, starting from listening to a set of low everyday sounds in a quiet environment, getting accustomed to one's own speech and other people's speech, getting accustomed to speech among background noise, etc.
History
The first hearing aids were ear trumpets, and were created in the 17th century. Some of the first hearing aids were external hearing aids. External hearing aids directed sounds in front of the ear and blocked all other noises. The apparatus would fit behind or in the ear.
The movement toward modern hearing aids began with the creation of the telephone, and the first electric hearing aid, the "akouphone", was created about 1895 by Miller Reese Hutchison. By the late 20th century, digital hearing aids were commercially available.
The invention of the carbon microphone, transmitters, digital signal processing chip or DSP, and the development of computer technology helped transform the hearing aid to its present form.
History of digital aids
The history of DHA can be divided into three stages. The first stage began in the 1960s with the widespread use of digital computers for simulation of audio processing and for the analysis of systems and algorithms. The work was conducted with the help of the very large digital computers of that era. These efforts were not actual digital hearing aids because the computers were not fast enough for audio processing in real time and their size prevented them from being described as wearable, but they allowed successful studies of the various hardware circuits and algorithms for digital processing of audio signals. The software package Block of Compiled Diagrams (BLODI) developed by Kelly, Lockbaum and Vysotskiy in 1961 allowed simulation of any sound system that could be characterized in the form of a block diagram. A special phone was created so that a person with a hearing impairment could listen to the digitally processed signals, but not in real time. In 1967, Harry Levitt used BLODI to simulate a hearing aid on a digital computer.
Almost ten years later the second stage began with the creation of the hybrid hearing aid, in which the analog components of a conventional hearing aid consisting of amplifiers, filters and signal limiting were combined with a separate digital programmable component in a conventional hearing aid case. The audio processing remained analog but it was controlled by the digital programmable component. The digital component could be programmed by connecting the device to an external computer in the laboratory then disconnected to allow the hybrid device to function as a conventional wearable hearing aid.
The hybrid device was effective from a practical point of view because of the low power consumption and compact size. At that time, low-power analog amplifier technology was well developed in contrast to the available semiconductor chips able to process digital audio in real time. The combination of high performance analog components for real time audio processing and a separate low power digital programmable component only for controlling the analog signal led to the creation of several low power digital programmable components able to implement different types of control.
A hybrid hearing aid was developed by Etymotic Design. A little later, Mangold and Lane created a programmable multi-channel hybrid hearing aid. Graupe with co-authors developed a digital programmable component that implemented an adaptive noise filter.
The third stage began in the early 1980s by a research group at Central Institute for the Deaf headed up by faculty members at Washington University in St. Louis MO. This group created the first fully digital wearable hearing aid. They first conceived a complete, comprehensive full digital hearing aid, then designed and fabricated, miniaturized full digital computer chips using custom digital signal processing chips with low power and very large scale integrated (VLSI) chip technology able to process both the audio signal in real time and the control signals, yet able to be powered by a battery and be fully wearable as a full digital wearable hearing aid able to be actually used by individuals with hearing loss in real-world environments. Engebretson, Morley and Popelka were the inventors of the first full digital hearing aid. Their work resulted in US Patent 4,548,082, "Hearing aids, signal supplying apparatus, systems for compensating hearing deficiencies, and methods" by A Maynard Engebretson, Robert E Morley Jr. and Gerald R Popelka, filed in 1984 and issued in 1985. This full digital wearable hearing aid also included many additional features now used in all contemporary full digital hearing aids including a bidirectional interface with an external computer, self-calibration, self-adjustment, wide bandwidth, digital programmability, a fitting algorithm based on audibility, internal storage of digital programs, and fully digital multichannel amplitude compression and output limiting. This group created several of these full digital hearing aids and used them for research on hearing impaired people as they wore them in the same manner as conventional hearing aids in real-world situations. In this first full DHA all stages of sound processing and control were carried out in binary form. The external sound was picked up by a microphone positioned in an ITE ear module to take advantage of the acoustic effects of the pinna, then converted into binary code, digitally processed and digitally controlled in real time, then converted back to an analog signal sent to two miniature loudspeakers positioned in the same ITE ear module. The ITE module also contained an inward facing microphone to measure the sound actually generated in the ear canal, a precursor to separate probe tube measures now routinely used for hearing aid fitting. The necessary electronic components, including batteries, to support this arrangement were situated in a BTE module that could be supplemented with a body worn module. These specialized hearing aid chips continued to become smaller, increase in computational ability and require even less power. Now, virtually all commercial hearing aids are fully digital and their digital signal processing capability has significantly increased. Very small and very low power specialized digital hearing aid chips are now used in all hearing aids manufactured worldwide. Many additional new features also have been added with various on-board advanced wireless technology.
Regulation
Canada
Hearing aids are Class II regulated medical devices under Canada's Food and Drugs Act.
Under Health Canada, the Medical Devices Directorate (MDD) regulates the safety, quality, and effectiveness of hearing aids. All hearing aids imported and sold in Canada are subject to a pre-market review. Post-market, Health Canada monitors the performance of the hearing aid and any consumer complaints.
Hearing aid financial assistance is available at both the federal and provincial level. Provincial hearing aid assistance and coverage can vary widely depending on the province and territory.
In Canada, a prescription is required to purchase hearing aids. Only licensed audiologists, Ear, Nose and Throat (ENT) doctors, hearing instrument practitioners (where the profession exists), and audioprothésistes (in Quebec) can prescribe hearing aids. Over-the-counter (OTC) hearing aids are currently not available for sale in Canada.
Canadian taxpayers can claim tax relief for hearing aids as a medical expense.
Ireland
Like much of the Irish health care system, hearing aid provision is a mixture of public and private.
Hearing aids are provided by the state to children, OAPs and to people whose income is at or below that of the state pension. The Irish state hearing aid provision is extremely poor; people often have to wait for two years for an appointment.
It is estimated that the total cost to the state of supplying one hearing aid exceeds €2,000.
Hearing aids are also available privately, and there is grant assistance available for insured workers. For the fiscal year ending 2016, the grant stands at a maximum of €500 per ear.
Irish taxpayers can also claim tax relief at the standard rate as hearing aids are recognised as a medical device.
Hearing aids in the Republic of Ireland are exempt from VAT.
Hearing aid providers in Ireland mostly belong to the Irish Society of Hearing Aid Audiologists.
United States
Ordinary hearing aids are Class I regulated medical devices under Federal Food and Drug Administration (FDA) rules. A 1976 statute explicitly prohibits any state requirement that is "different from, or in addition to, any requirement applicable" to regulated medical devices (which includes hearing aids) which relates "to the safety and effectiveness of the device". Inconsistent state regulation is preempted under the federal law. In the late 1970s, the FDA established federal rules governing hearing aid sales, and addressed various requests by state authorities for exemptions from federal preemption, granting some and denying others. The Over-the-Counter Hearing Aid Act (OTC Act) was passed under the FDA Reauthorization Act of 2017, creating a class of hearing aids regulated by the FDA available directly to consumers without involvement from a licensed professional. This law's provisions are expected to go into effect in 2020.
In August 2022, the FDA issued a final rule to improve access to hearing aids. The action establishes a new category of over-the-counter (OTC) hearing aids, enabling consumers with perceived mild to moderate hearing impairment to purchase hearing aids directly from stores or online retailers without the need for a medical exam, prescription or a fitting adjustment by an audiologist. The FDA action amends existing rules that apply to prescription hearing aids for consistency with the new OTC category, it repeals the conditions for sale for hearing aids, and it includes provisions that address some of the effects of the FDA OTC hearing aid regulations on state regulation of hearing aids. The FDA also issued the final guidance, Regulatory Requirements for Hearing Aid Devices and Personal Sound Amplification Products (PSAPs), to clarify the differences between hearing aids, which are medical devices, and PSAPs, consumer products that help people with normal hearing amplify sounds.
Cost
Several industrialized countries supply free or heavily discounted hearing aids through their publicly funded health care system.
Australia
The Australian Department of Health and Ageing provides eligible Australian citizens and residents with a basic hearing aid free-of-charge, though recipients can pay a "top up" charge if they wish to upgrade to a hearing aid with more or better features. Maintenance of these hearing aids and a regular supply of batteries is also provided, on payment of a small annual maintenance fee.
Canada
In Canada, health care is a responsibility of the provinces. In the province of Ontario, the price of hearing aids is partially reimbursed through the Assistive Devices Program of the Ministry of Health and Long-Term care, up to $500 for each hearing aid. Like eye appointments, audiological appointments are no longer covered through the provincial public health plan. Audiometric testing can still easily be obtained, often free of charge, in private sector hearing aid clinics and some ear, nose and throat doctors offices. Hearing aids may be covered to some extent by private insurance or in some cases through government programs such as Veterans Affairs Canada or Workplace Safety & Insurance Board.
Iceland
Social Insurance pays a one time fee of ISK 30,000 for any kind of hearing aid. However, the rules are complicated and require that both ears have significant hearing loss in order to qualify for reimbursement. BTE hearing aids range from ISK 60,000 to ISK 300,000.
India
In India hearing aids of all kinds are easily available. Under central and state government health services, the poor can often avail themselves of free hearing devices. However, market prices vary for others and can range from Rs 10,000 to Rs 275,000 per ear.
United Kingdom
From 2000 to 2005 the Department of Health worked with Action on Hearing Loss (then called RNID) to improve the quality of NHS hearing aids so every NHS audiology department in England was fitting digital hearing aids by March 2005. By 2003 over 175,000 NHS digital hearing aids had been fitted to 125,000 people. Private companies were recruited to enhance the capacity, and two were appointed – David Ormerod Hearing Centres, partly owned by Alliance Boots and Ultravox Group, a subsidiary of Amplifon.
Within the UK, the NHS provides digital BTE hearing aids to NHS patients, on long-term loan, free of charge. Other than BAHAs (bone anchored hearing aid) or cochlear implants, where specifically required, BTEs are usually the only style available. Private purchases may be necessary if a user desires a different style. Batteries are free.
In 2014 the Clinical Commissioning Group in North Staffordshire considered proposals to end provision of free hearing aids for adults with mild to moderate age related hearing loss, which currently cost them £1.2m a year. Action on Hearing Loss mobilised a campaign against the proposal.
In June 2018 the National Institute for Health and Care Excellence produced new guidance saying that hearing aids should be offered at the first opportunity when hearing loss affects the individual's ability to hear and communicate, rather than waiting for arbitrary thresholds of hearing loss to be reached.
United States
Most private healthcare providers in the United States do not provide coverage for hearing aids, so all costs are usually borne by the recipient. The cost for a single hearing aid can vary between $500 and $6,000 or more, depending on the level of technology and whether the clinician bundles fitting fees into the cost of the hearing aid. Though if an adult has hearing loss which substantially limits major life activities, some state-run vocational rehabilitation programs can provide upwards of full financial assistance. Severe and profound hearing loss often falls within the "substantially limiting" category. Less expensive hearing aids can be found on the internet or mail order catalogs, but most in the under-$200 range tend to amplify the low frequencies of background noise, making it harder to hear the human voice.
Military veterans receiving VA medical care are eligible for hearing aids based on medical need. The Veterans Administration pays the full cost of testing and hearing aids to qualified military veterans. Major VA medical facilities provide complete diagnostic and audiology services.
The cost of hearing aids is a tax-deductible medical expense for those who itemize medical deductions.
Research involving more than 40,000 US households showed a convincing correlation between the degree of hearing loss and the reduction of personal income. According to the same research, hearing aids were shown to mitigate the impact of income loss by 90%–100% for those with milder hearing losses and from 65%–77% for those with severe to moderate hearing loss.
Batteries
While there are some instances that a hearing aid uses a rechargeable battery or a long-life disposable battery, the majority of modern hearing aids use one of five standard button cell zinc–air batteries. (Older hearing aids often used mercury battery cells, but these cells have become banned in most countries today.) Modern hearing aid button cell types are typically referred to by their common number name or the color of their packaging.
They are typically loaded into the hearing aid via a rotating battery door, with the flat side (case) as the positive terminal (cathode) and the rounded side as the negative terminal (anode).
These batteries all operate from 1.35 to 1.45 volts.
The type of battery a specific hearing aid utilizes depends on the physical size allowable and the desired lifetime of the battery, which is in turn determined by the power draw of the hearing aid device. Typical battery lifetimes run between 1 and 14 days (assuming 16-hour days).
| Technology | Devices | null |
520449 | https://en.wikipedia.org/wiki/Pacific%20plate | Pacific plate | The Pacific plate is an oceanic tectonic plate that lies beneath the Pacific Ocean. At , it is the largest tectonic plate.
The plate first came into existence as a microplate 190 million years ago, at the triple junction between the Farallon, Phoenix, and Izanagi plates. The Pacific plate subsequently grew to where it underlies most of the Pacific Ocean basin. This reduced the Farallon plate to a few remnants along the west coast of the Americas and the Phoenix plate to a small remnant near the Drake Passage, and destroyed the Izanagi plate by subduction under Asia.
The Pacific plate contains an interior hot spot forming the Hawaiian Islands.
Boundaries
The north-eastern side is a divergent boundary with the Explorer plate, the Juan de Fuca plate and the Gorda plate forming respectively the Explorer Ridge, the Juan de Fuca Ridge and the Gorda Ridge. In the middle of the eastern side is a transform boundary with the North American plate along the San Andreas Fault, and a divergent boundary with the Cocos plate. The south-eastern side is a divergent boundary with the Nazca plate forming the East Pacific Rise.
The southern side is a divergent boundary with the Antarctic plate forming the Pacific–Antarctic Ridge.
The western side is bounded by the Okhotsk microplate at the Kuril–Kamchatka Trench and the Japan Trench. The plate forms a convergent boundary by subducting under the Philippine Sea plate creating the Mariana Trench, has a transform boundary with the Caroline plate, and has a collision boundary with the North Bismarck plate.
In the south-west, the Pacific plate has a complex but generally convergent boundary with the Indo-Australian plate, subducting under it north of New Zealand forming the Tonga Trench and the Kermadec Trench. The Alpine Fault marks a transform boundary between the two plates, and further south the Indo-Australian plate subducts under the Pacific plate forming the Puysegur Trench. The southern part of Zealandia, which is to the east of this boundary, is the plate's largest block of continental crust. Hillis and Müller are reported to consider the Bird's Head plate to be moving in unison with the Pacific plate, but Bird considers them to be unconnected.
The northern side is a convergent boundary subducting under the North American plate forming the Aleutian Trench and the corresponding Aleutian Islands (see also: Aleutian Arc).
Paleo-geology of the Pacific plate
The Pacific plate is almost entirely oceanic crust, but it contains some continental crust in New Zealand, Baja California, and coastal California.
The Pacific plate has the distinction of showing one of the largest areal sections of the oldest members of seabed geology being entrenched into eastern Asian oceanic trenches. A geologic map of the Pacific Ocean seabed shows not only the geologic sequences, and associated Ring of Fire zones on the ocean's perimeters, but the various ages of the seafloor in a stairstep fashion, youngest to oldest, the oldest being consumed into the Asian oceanic trenches. The oldest part disappearing by way of the plate tectonics cycle is early-Cretaceous (145 to 137 million years ago).
The Pacific plate originated at the triple junction of the three main oceanic plates of Panthalassa, the Farallon, Phoenix, and Izanagi plates, around 190 million years ago. The plate formed because the triple junction had converted to an unstable form surrounded on all sides by transform faults, due to the development of a kink in one of the plate boundaries. The "Pacific Triangle", the oldest part of the Pacific plate, created during the initial stages of plate formation, is located just east of the Mariana Trench. The growth of the Pacific plate reduced the Farallon plate to a few remnants along the west coast of the Americas (such as the Juan de Fuca plate) and the Phoenix plate to a small remnant near the Drake Passage, and destroyed the Izanagi plate by subduction under Asia.
| Physical sciences | Tectonic plates | Earth science |
520635 | https://en.wikipedia.org/wiki/Solifugae | Solifugae | Solifugae is an order of arachnids known variously as solifuges, sun spiders, camel spiders, and wind scorpions. The order includes more than 1,000 described species in about 147 genera. Despite the common names, they are neither true scorpions (order Scorpiones) nor true spiders (order Araneae). Because of this, it is less ambiguous to call them "solifuges". Most species of solifuge live in dry climates and feed opportunistically on ground-dwelling arthropods and other small animals. The largest species grow to a length of , including legs. A number of urban legends exaggerate the size and speed of solifuges, and their potential danger to humans, which is negligible.
Etymology
The order's name is derived from the Latin "sol" meaning "sun" and "fugere" meaning "to flee". Put together, it means "those who flee from the sun". These animals have a number of common names including sun spiders, wind scorpions, wind spiders, red romans, and camel spiders. In Afrikaans, they are known as "haarskeerders" ("hair cutters"), and "baardskeerders" ("beard cutters"). This is in reference to myths that they cut hair to be used as nest bedding.
Anatomy and physiology
Solifuges are moderately small to large arachnids (a few millimeters to several centimeters in body length), with the larger species reaching in length, including legs. In practice, the respective lengths of the legs of various species differ greatly, so the resulting figures are often misleading. More practical measurements refer primarily to the body length, quoting leg lengths separately, if at all. The body length is up to . Most species are closer to long, and some small species are under in head-plus-body length when mature.
Like that of spiders, the body plan of the Solifugae has two main tagmata: the prosoma, or cephalothorax, is the anterior tagma, and the 10-segmented abdomen, or opisthosoma, is the posterior tagma. The abdominal tergites and sternites are separated by large areas of intersegmental membranes, giving it a high degree of flexibility and ability to stretch considerably, which allows it to consume a large amount of food. As shown in the illustrations, the solifuge prosoma and opisthosoma are not separated by nearly as clear a constriction and connecting tube or "pedicel" as occurs in Araneae. The lack of the pedicel reflects another difference between the Solifugae and spiders, namely that solifuges lack both spinnerets and silk, and do not spin webs. Spiders need considerable mobility of their abdomens in their spinning activities, and the Solifugae have no such adaptation.
The prosoma comprises the head, the mouthparts, and the somites that bear the legs and the pedipalps. It is covered by a carapace, also called a prosomal dorsal shield or peltidium, which is composed of three distinct elements called propeltidium, mesopeltidium and metapeltidium. The propeltidium contains the eyes, the chelicerae that, in most species, are conspicuously large, the pedipalps and the first two pairs of legs. Meso- and metapeltidium contains the third and fourth pairs of legs. The chelicerae serve as jaws and in many species also are used for stridulation. Unlike scorpions, solifuges do not have a third tagma that forms a "tail".
Currently, neither fossil nor embryological evidence shows that arachnids ever had a separate thorax-like division, so the validity of the term cephalothorax, which means a fused cephalon, or head, and thorax, has been questioned. Also, arguments exist against use of "abdomen", as the opisthosoma of many arachnids contains organs atypical of an abdomen, such as a heart and respiratory organs.
Like other arachnids outside the orders of scorpions and the Tetrapulmonata, the Solifugae lack book lungs, having instead a well-developed tracheal system that inhales and exhales air through a number of spiracles - one pair between the second and third pair of walking legs, two pairs on the abdomen on abdominal segments three and four, and an unpaired spiracle on the fifth abdominal segment. Air sacs are attached to the branching tracheae, with tracheoles penetrating the epithelia of internal organs. Hemocyanin, a respiratory pigment common in the hemolymph of many arachnids and other arthropods, is absent. As embryos they also have opisthosomal protuberances resembling the pulmonary sacs found in some palpigrades.
Chelicerae
Among the most distinctive features of the Solifugae are their large chelicerae, which in many species are longer than the prosoma. Each of the two chelicerae has two articles (segments, parts connected by a joint), forming a powerful pincer, much like that of a crab; each article bears a variable number of teeth, largely depending on the species. The chelicerae of many species are surprisingly strong; they are capable of shearing hair or feathers from vertebrate prey or carrion, and of cutting through skin and thin bones such as those of small birds. Many Solifugae stridulate with their chelicerae, producing a rattling noise.
Legs and pedipalps
These elements work the same way as in most other arachnids. Although the Solifugae appear to have five pairs of legs, only the hind four pairs are true legs. Each true leg has seven segments: coxa, trochanter, femur, patella, tibia, metatarsus, and tarsus.
The first, or anterior, of the five pairs of leg-like appendages are not "actual" legs, but pedipalps, and they have only five segments each. The pedipalps of the Solifugae function partly as sense organs similar to insects' antennae, and partly in locomotion, feeding, and fighting. In normal locomotion, they do not quite touch the ground, but are held out to detect obstacles and prey; in that attitude, they look particularly like an extra pair of legs or perhaps arms. Reflecting the great dependence of the Solifugae on their tactile senses, their anterior true legs commonly are smaller and thinner than the posterior three pairs. That smaller anterior pair acts largely in a sensory role as a supplement to the pedipalps, and in many species they accordingly lack tarsi. At the tips of their pedipalps, Solifugae bear a membranous suctorial organ, which are used for capturing prey, and also for bringing water to their mouthparts for drinking and climbing smooth surfaces.
For the most part, only the posterior three pairs of legs are used for running. On the undersides of the coxae and trochanters of the last pair of legs, the Solifugae have fan-shaped sensory organs called malleoli or racquet (or racket) organs. Sometimes, the blades of the malleoli are directed forward, sometimes not. They have been suspected to be sensory organs for the detection of vibrations in the soil, perhaps to detect threats and potential prey or mates. These structures may be chemoreceptors.
Males are usually smaller than females, with relatively longer legs. Unlike females, the males bear a pair of flagella, one on each chelicera. In the accompanying photograph of a male solifuge, one flagellum is just visible near the tip of each chelicera. The flagella, which bend back over the chelicerae, are sometimes called horns and are believed to have some sexual connection, but their function has not yet been clearly explained.
Eyes
Solifuges have a pair of large central eyes known as median ocelli These eyes are oriented at the very front of its cephalothorax and are placed close together. These eyes have a pigment-cup structure and are covered by a domed outer lens made from the animal's exoskeleton. Below the dome is the animal's retina, a multi-tiered structure with a layer of cells called the vitreous body at its top. Underneath is the thin preretinal membrane, acting as a barrier between the vitreous body above and the rhabdomeres beneath. Rhabdomeres are light-sensitive and function as the eye's photoreceptors. Interspersed between the rhabdomeres are pigment cells. The eye's optic nerve begins at its center and is connected to the axons of numerous rhabdomeres.
In addition to the median eyes, solifuges possess a pair of vestigial lateral ocelli. These eyes are found in pits on the animal's cephalic lobes near the chelicerae. The ocelli's lenses are usually atrophied. However, in some species both nerves and pigment cells are present. In species where lateral eyes are functional, they probably aid in detecting motions or changes in light intensity.
Habitat and distribution
Most solifuges live in tropics and subtropical deserts in the Americas, Southern Europe, Africa, the Middle East, and South Asia. Surprisingly, these animals are absent in Australia and Madagascar. Within the desert, solifuges live in a variety of micro-habitats. These include sand dunes, sand flats, floodplains, rocky hillsides, desert shrublands, gravel plains, and mountain valleys. In addition to the desert, certain solifuges live in more arid grasslands and forests.
Depending on the species in question, solifuges may be more sedentary or on the move. Sedentary species are often fossorial, living in relatively permanent burrows underground. Transitory species spend most of their time up the surface, occasionally seeking refuge in cracks or under rocks and vegetation.
Behavior and life history
Diet and hunting
Solifuges are carnivores and typically generalists, feeding on a wide variety of prey in their given environment. For most species, insects make up the bulk of their diet. However, these animals have been known to consume anything they can subdue. This includes other arachnids like spiders, scorpions, and smaller solifuges, other arthropods like millipedes, and small lizards, birds, and mammals. Additionally, solifuges are voracious eaters. It's common for adult females to eat so much that they're temporarily unable to walk.
When looking for prey, most solifuges rapidly move about while tapping their pedipalps on the ground. The only exception is the majority of termite-loving species, as they prefer to be more sedentary. In addition to using their pedipalps, solifuges have a variety of methods to locate prey. These include seeing movements with their eyes, feeling with their long hairlike setae, smelling with their malleoli, and sensing vibrations. How much the animal relies on each sense depends on the species. While all hunt on the ground, some species are great climbers, able to search for prey on trees, shrubs, and on artificial structures.
Solifuges hunt their prey using three main hunting strategies: stalking, chasing, and ambushing. Depending on the meal's size, prey is seized with the animal's pedipalps or massive chelicerae. When the pedipalps are used, prey is initially caught with the limb's suction cups, then rapidly pulled towards the chelicerae to be chewed. These motions happen so fast that they can't be distinguished. Before eating, solifuges prepare their food by removing any parts they find unfavorable. In arthropods, these are typically areas that have a high amount of chitin (heads, antennae, wings, etc).
Solifuges eat in different ways based on the shape of their food. Prey that is long and narrow is held perpendicular to the chelicerae and chewed from one end to another. More round prey is chewed by rotating the body all at once. This chewing motion turns the food into a liquidized paste which is then swallowed by the animal's pharynx. Solifuges that haven't fed for long periods are known to eat faster than ones that fed recently. Larger solifuges are also known to eat faster than smaller ones.
Reproduction
The Solifugae are typically univoltine (reproducing once a year). Reproduction can involve direct or indirect sperm transfer; when indirect, the male emits a spermatophore on the ground and then inserts it with his chelicerae in the female's genital pore. To do this, he flings the female on her back.
The female then digs a burrow, into which she lays 50 to 200 eggs; some species then guard them until they hatch. Because the female does not feed during this time, she tries to fatten herself beforehand, and a species of has been observed to eat more than 100 flies during that time in the laboratory. The Solifugae undergo a number of stages including, egg, postembryo, 9–10 nymphal instars, and adults.
Classification and phylogeny
Solifuges are an order of arachnids comprising over 1200 species in 146 genera assigned to 16 different families. Solifuges can be divided into two groups of families which are recognized as distinct suborders. These are the Australosolifugae which live predominantly in the Southern Hemisphere and the Boreosolifugae which live mostly in the Northern Hemisphere. This phylogeny is considered congruent with a Gondwanan origin for Australosolifugae and a Laurasian origin for Boreosolifugae. When looking at their relationships, the families Ammotrechidae and Daesiidae were found to be paraphyletic, leading to multiple clades without a name. Because of this, a later genomic study established three additional families: Dinorhaxidae, Lipophagidae, and Namibesiidae.
Suborder Boreosolifugae
Eremobatidae Kraepelin, 1901
Galeodidae Sundevall, 1833
Gylippidae Roewer, 1933
Karschiidae Kraepelin, 1899
Rhagodidae Pocock, 1897
Suborder Australosolifugae
Ammotrechidae Roewer, 1934
Ceromidae Roewer, 1933
Daesiidae Kraepelin, 1899
Dinorhaxidae (Roewer, 1933)
Hexisopodidae Pocock, 1897
Melanoblossiidae Roewer, 1933
Mummuciidae Roewer, 1934
Lipophagidae (Wharton, 1981)
Namibesiidae (Wharton, 1981)
Solpugidae Leach, 1815
incertae sedis
†Protosolpugidae Petrunkevitch, 1953
Phylogeny
Below is a family tree of the various solifuge families based on phylogenomics.
Relationship with humans
Solifuges have been recognized as distinct taxa from ancient times. In Aelian's De natura animalium, "four-jawed spiders" are credited, along with scorpions, as being responsible for the abandoning of a desert region near the Astaboras river (said to be in India, but thought to be a river in Ethiopia). Anton August Heinrich Lichtenstein theorized in 1797 that the "mice" that plagued the Philistines in the Old Testament were Solifugae. During World War I, troops stationed in Abū Qīr, Egypt, would stage fights between captive "jerrymanders", as they referred to them, and placed bets on the outcome. Similarly, British troops stationed in Libya in World War II staged fights between solifuges and scorpions.
Urban legends
The Solifugae are the subject of many legends and exaggerations about their size, speed, behavior, appetite, and lethality. They are not especially large, the biggest having a leg span around . They are fast on land compared to other invertebrates, with their top speed estimated to be .
The Solifugae apparently have neither venom glands nor any venom-delivery apparatus such as the fangs of spiders, stings of wasps, or venomous setae of caterpillars (e.g., Lonomia or Acharia species). One 1978 study is frequently quoted, in which the authors report detection of an exception in India, in that Rhagodes nigrocinctus had venom glands, and that injection of the secretion into mice was frequently fatal. However, no supporting studies have confirmed either statement, such as by independent detection of the glands as claimed, or the relevance of the observations, if correct. Even the authors of the original account admitted to having found no means of delivery of the putative venom by the animal, and the only means of administering the material to the mice was by parenteral injection. Given that many non-venoms such as saliva, blood and glandular secretions can be lethal if injected, and that no venomous function was even speculated upon in this study, there is still no evidence for even one venomous species of solifuge.
Because of their unfamiliar spider-like appearance and rapid movements, Solifugae have startled or even frightened many people. This fear was sufficient to drive a family from their home in August 2008 when one was allegedly discovered in a soldier's house in Colchester, England, and caused the family to blame the solifuge for the death of their pet dog. An Arizona resident developed painful lesions due to a claimed solifuge bite but could not produce a specimen for confirmation. Though they are not venomous, the powerful chelicerae of a large specimen may inflict a painful nip, but nothing medically significant.
Claims that Solifugae aggressively chase people are also untrue, as they are merely trying to stay in the shade/shadow provided by the human.
| Biology and health sciences | Arachnids | Animals |
521047 | https://en.wikipedia.org/wiki/Cocos%20plate | Cocos plate | The Cocos plate is a young oceanic tectonic plate beneath the Pacific Ocean off the west coast of Central America, named for Cocos Island, which rides upon it. The Cocos plate was created approximately 23 million years ago when the Farallon plate broke into two pieces, which also created the Nazca plate. The Cocos plate also broke into two pieces, creating the small Rivera plate. The Cocos plate is bounded to the northeast by the North American plate and the Caribbean plate. To the west it is bounded by the Pacific plate and to the south by the Nazca plate.
The only land above water on the Cocos plate is Cocos Island, which is administered by Costa Rica and lies approximately 550 km (342 mi; 297 nmi) southwest of the Costa Rican mainland.
Geology
The Cocos plate was created by sea floor spreading along the East Pacific Rise and the Cocos Ridge, specifically in a complicated area geologists call the Cocos-Nazca spreading system. From the rise the plate is pushed eastward and pushed or dragged (perhaps both) under the less dense Caribbean plate, in the process called subduction. The subducted leading edge heats up and adds its water to the mantle above it. In the mantle layer called the asthenosphere, mantle rock melts to make magma, trapping superheated water under great pressure. As a result, to the northeast of the subducting edge lies the continuous arc of volcanos – also known as the Central America Volcanic Arc – stretching from Costa Rica to Guatemala, and a belt of earthquakes that extends farther north, into Mexico.
The northern boundary of the Cocos plate is the Middle America Trench. The eastern boundary is a transform fault, the Panama fracture zone. The southern boundary is a mid-oceanic ridge, the Cocos–Nazca spreading centre. The western boundary is another mid-ocean ridge, the East Pacific Rise.
A hotspot under the Galápagos Islands lies along the Galápagos Rise. (see Galápagos hotspot and Galápagos microplate)
The Rivera plate, north of the Cocos plate, is thought to have separated from the Cocos plate 5–10 million years ago. The boundary between the two plates appears to lack a definite transform fault, yet they are regarded as distinct. After its separation from the Cocos plate, the Rivera plate started acting as an independent microplate.
The devastating 1985 Mexico City earthquake and the 2017 Chiapas earthquake were results of the subduction of the Cocos plate beneath the North American plate. The devastating El Salvador earthquakes in January 2001 and February 2001 were generated by the subduction of this plate beneath the Caribbean plate.
| Physical sciences | Tectonic plates | Earth science |
521060 | https://en.wikipedia.org/wiki/Philippine%20Sea%20plate | Philippine Sea plate | The Philippine Sea plate or the Philippine plate is a tectonic plate comprising oceanic lithosphere that lies beneath the Philippine Sea, to the east of the Philippines. Most segments of the Philippines, including northern Luzon, are part of the Philippine Mobile Belt, which is geologically and tectonically separate from the Philippine Sea plate.
The plate is bordered mostly by convergent boundaries: To the north, the Philippine Sea plate meets the Okhotsk microplate at the Nankai Trough. The Philippine Sea plate, the Amurian plate, and the Okhotsk plate meet near Mount Fuji in Japan. The thickened crust of the Izu–Bonin–Mariana arc colliding with Japan constitutes the Izu Collision Zone. The east of the plate includes the Izu–Ogasawara (Bonin) and the Mariana Islands, forming the Izu–Bonin–Mariana Arc system. There is also a divergent boundary between the Philippine Sea plate and the small Mariana plate which carries the Mariana Islands. To the east, the Pacific plate subducts beneath the Philippine Sea plate at the Izu–Ogasawara Trench. To the south, the Philippine Sea plate is bounded by the Caroline plate and Bird's Head plate. To the west, the Philippine Sea plate subducts under the Philippine Mobile Belt at the Philippine Trench and the East Luzon Trench. (The adjacent rendition of Prof. Peter Bird's map is inaccurate in this respect.) To the northwest, the Philippine Sea plate meets Taiwan and the Nansei islands on the Okinawa plate, and southern Japan on the Amurian plate.
It also meets the Yangtze plate due northwest.
| Physical sciences | Tectonic plates | Earth science |
521081 | https://en.wikipedia.org/wiki/African%20plate | African plate | The African plate, also known as the Nubian plate, is a major tectonic plate that includes much of the continent of Africa (except for its easternmost part) and the adjacent oceanic crust to the west and south. It also includes a narrow strip of Western Asia along the Mediterranean Sea, including much of Israel and Lebanon. It is bounded by the North American plate and South American plate to the west (separated by the Mid-Atlantic Ridge); the Arabian plate and Somali plate to the east; the Eurasian plate, Aegean Sea plate and Anatolian plate to the north; and the Antarctic plate to the south.
Between and , the Somali plate began rifting from the African plate along the East African Rift. Since the continent of Africa consists of crust from both the African and the Somali plates, some literature refers to the African plate as the Nubian plate to distinguish it from the continent as a whole.
Boundaries
The western edge of the African plate is a divergent boundary with the North American plate to the north and the South American plate to the south which forms the central and southern part of the Mid-Atlantic Ridge. The African plate is bounded on the northeast by the Arabian plate, the southeast by the Somali plate, the north by the Eurasian plate, the Aegean Sea plate, and the Anatolian plate, and on the south by the Antarctic plate at the Southwest Indian Ridge. All of these are divergent or spreading boundaries with the exception of the northern boundary and a short segment near the Azores known as the Terceira Rift.
Components
The African plate includes several cratons, stable blocks of old crust with deep roots in the subcontinental lithospheric mantle, and less stable terranes, which came together to form the African continent during the assembly of the supercontinent Pangea around 250 million years ago. The cratons are from south to north, the Kalahari Craton, Congo Craton, Tanzania Craton and West African Craton. The cratons were widely separated in the past, but came together during the Pan-African orogeny and stayed together when Gondwana split up. The cratons are connected by orogenic belts, regions of highly deformed rock where the tectonic plates have engaged.
The Saharan Metacraton has been tentatively identified as the remains of a craton that has become detached from the subcontinental lithospheric mantle, but alternatively may consist of a collection of unrelated crustal fragments swept together during the Pan-African orogeny.
In some areas, the cratons are covered by sedimentary basins, such as the Tindouf Basin, Taoudeni Basin and Congo Basin, where the underlying archaic crust is overlaid by more recent Neoproterozoic sediments. The plate includes shear zones such as the Central African Shear Zone (CASZ) where, in the past, two sections of the crust were moving in opposite directions, and rifts such as the Anza Trough where the crust was pulled apart, and the resulting depression filled with more modern sediment.
Modern movements
The African plate is rifting in the eastern interior of the African continent along the East African Rift. This rift zone separates the African plate to the west from the Somali plate to the east. One hypothesis proposes a mantle plume rising beneath the Afar region pushing the crust outward, whereas an opposing hypothesis explains the rifting by dynamics in the crust, as a break in the African plate along a line of maximum weakness as plates to its east move rapidly northward.
The African plate's speed is estimated at per year. It has been moving over the past 100 million years or so in a general northeast direction. It is pushing closer to the Eurasian plate, causing subduction where oceanic crust is converging with continental crust (e.g. portions of the central and eastern Mediterranean). In the western Mediterranean, the relative motions of the Eurasian and African plates produce a combination of lateral and compressive forces, concentrated in a zone known as the Azores–Gibraltar Fault Zone. Along its northeast margin, the African plate is bounded by the Red Sea Rift where the Arabian plate is moving away from the African plate.
The New England hotspot in the Atlantic Ocean has probably created a short line of mid- to late-Tertiary age seamounts on the African plate but appears to be currently inactive.
| Physical sciences | Tectonic plates | Earth science |
521248 | https://en.wikipedia.org/wiki/Tropical%20forest | Tropical forest | Tropical forests are forested ecoregions with tropical climates – that is, land areas approximately bounded by the tropics of Cancer and Capricorn, but possibly affected by other factors such as prevailing winds.
Some tropical forest types are difficult to categorize. While forests in temperate areas are readily categorized on the basis of tree canopy density, such schemes do not work well in tropical forests. There is no single scheme that defines what a forest is, in tropical regions or elsewhere. Because of these difficulties, information on the extent of tropical forests varies between sources. However, tropical forests are extensive, making up just under half the world's forests. The tropical domain has the largest proportion of the world's forests (45 percent), followed by the boreal, temperate and subtropical domains.
More than 3.6 million hectares of virgin tropical forest was lost in 2018.
History
The original tropical rainforests, which covered the planet's land surface, were the type of flora that covered Earth.Other canopy forests expanded north-south of the equator during the Paleogene epoch, around 40 million years ago, as a result of the emergence of drier, cooler climates.
The tropical forest was originally identified as a specific type of biome in 1949.
Types of tropical forest
Tropical forests are often thought of as evergreen rainforests and moist forests, but these account for only a portion of them (depending on how they are defined – see maps). The remaining tropical forests are a diversity of many different forest types including:
Eucalyptus open forest, tropical coniferous forests, savanna woodland (e.g. Sahelian forest), and mountain forests (the higher elevations of which are cloud forests). Over even relatively short distances, the boundaries between these biomes may be unclear, with ecotones between the main types.
The nature of tropical forests in any given area is affected by several factors, most importantly:
Geographical: location and climatic zone (see sub-types), with:
Temperature profile, which is relatively even in equatorial rainforest or with a cooler season towards subtropical latitudes;
Precipitation levels and seasonality, with strong dry seasons significantly affecting flora (e.g. the predominance of lianas);
Elevation affects the above, often creating "ecological islands" with high endemism (e.g. Mount Kinabalu in the Borneo rainforest).
Historical: prehistoric age of forest and level of recent disturbance (see threats), changing primary (usually maximum biodiversity) into secondary forest, degenerating into bamboo forest after prolonged swidden agriculture (e.g. in several areas of Indo-China).
Soil characteristics (also subject to various classifications): including depth and drainage.
The Global 200 scheme
The Global 200 scheme, promoted by the World Wildlife Fund, classifies three main tropical forest habitat types (biomes), grouping together tropical and sub-tropical areas (maps below):
Tropical and subtropical coniferous forests.
Tropical and subtropical dry broadleaf forests,
Tropical and subtropical moist broadleaf forests,
Extent of tropical and sub-tropical -
Threats
A number of tropical forests have been designated High-Biodiversity Wilderness Areas, but remain subject to a wide range of disturbances, including more localized pressures such as habitat loss and degradation and anthropogenic climate change.
Studies have also shown that ongoing climate change is increasing the frequency and intensity of some climate extremes (e.g. droughts, heatwaves and hurricanes) which, in combination with other local human disturbances, are driving unprecedented negative ecological consequences for tropical forests around the world. All tropical forests have experienced at least some levels of disturbance.
Current deforestation in the biodiversity hotspots of North of South America, sub-Saharan Africa, South-East Asia and the Pacific, can be attributed to export of commodities such as: beef, soy, coffee, cacao, palm oil, and timber; there is a requirement for "strong transnational efforts ... by improving supply chain transparency [and] public–private engagement".
A study in Borneo describes how, between 1973 and 2018, the old-growth forest had been reduced from 76% to 50% of the island, mostly due to fire and agricultural expansion. A widely-held view is that placing a value on the ecosystem services these forests provide may bring about more sustainable policies. However, clear monitoring and evaluation mechanisms for environmental, social and economic outcomes are needed.
For example, a study in Vietnam indicated that poor and inconsistent data combined with a lack of human resources and political interest (thus lack of financial support) are hampering efforts to improve forest land allocation and a Payments for Forest Environmental Services scheme.
| Physical sciences | Forests | Earth science |
521267 | https://en.wikipedia.org/wiki/Reflection%20%28physics%29 | Reflection (physics) | Reflection is the change in direction of a wavefront at an interface between two different media so that the wavefront returns into the medium from which it originated. Common examples include the reflection of light, sound and water waves. The law of reflection says that for specular reflection (for example at a mirror) the angle at which the wave is incident on the surface equals the angle at which it is reflected.
In acoustics, reflection causes echoes and is used in sonar. In geology, it is important in the study of seismic waves. Reflection is observed with surface waves in bodies of water. Reflection is observed with many types of electromagnetic wave, besides visible light. Reflection of VHF and higher frequencies is important for radio transmission and for radar. Even hard X-rays and gamma rays can be reflected at shallow angles with special "grazing" mirrors.
Reflection of light
Reflection of light is either specular (mirror-like) or diffuse (retaining the energy, but losing the image) depending on the nature of the interface. In specular reflection the phase of the reflected waves depends on the choice of the origin of coordinates, but the relative phase between s and p (TE and TM) polarizations is fixed by the properties of the media and of the interface between them.
A mirror provides the most common model for specular light reflection, and typically consists of a glass sheet with a metallic coating where the significant reflection occurs. Reflection is enhanced in metals by suppression of wave propagation beyond their skin depths. Reflection also occurs at the surface of transparent media, such as water or glass.
In the diagram, a light ray PO strikes a vertical mirror at point O, and the reflected ray is OQ. By projecting an imaginary line through point O perpendicular to the mirror, known as the normal, we can measure the angle of incidence, θi and the angle of reflection, θr. The law of reflection states that θi = θr, or in other words, the angle of incidence equals the angle of reflection.
In fact, reflection of light may occur whenever light travels from a medium of a given refractive index into a medium with a different refractive index. In the most general case, a certain fraction of the light is reflected from the interface, and the remainder is refracted. Solving Maxwell's equations for a light ray striking a boundary allows the derivation of the Fresnel equations, which can be used to predict how much of the light is reflected, and how much is refracted in a given situation. This is analogous to the way impedance mismatch in an electric circuit causes reflection of signals. Total internal reflection of light from a denser medium occurs if the angle of incidence is greater than the critical angle.
Total internal reflection is used as a means of focusing waves that cannot effectively be reflected by common means. X-ray telescopes are constructed by creating a converging "tunnel" for the waves. As the waves interact at low angle with the surface of this tunnel they are reflected toward the focus point (or toward another interaction with the tunnel surface, eventually being directed to the detector at the focus). A conventional reflector would be useless as the X-rays would simply pass through the intended reflector.
When light reflects off a material with higher refractive index than the medium in which is traveling, it undergoes a 180° phase shift. In contrast, when light reflects off a material with lower refractive index the reflected light is in phase with the incident light. This is an important principle in the field of thin-film optics.
Specular reflection forms images. Reflection from a flat surface forms a mirror image, which appears to be reversed from left to right because we compare the image we see to what we would see if we were rotated into the position of the image. Specular reflection at a curved surface forms an image which may be magnified or demagnified; curved mirrors have optical power. Such mirrors may have surfaces that are spherical or parabolic.
Laws of reflection
If the reflecting surface is very smooth, the reflection of light that occurs is called specular or regular reflection. The laws of reflection are as follows:
The incident ray, the reflected ray and the normal to the reflection surface at the point of the incidence lie in the same plane.
The angle which the incident ray makes with the normal is equal to the angle which the reflected ray makes to the same normal.
The reflected ray and the incident ray are on the opposite sides of the normal.
These three laws can all be derived from the Fresnel equations.
Mechanism
In classical electrodynamics, light is considered as an electromagnetic wave, which is described by Maxwell's equations. Light waves incident on a material induce small oscillations of polarisation in the individual atoms (or oscillation of electrons, in metals), causing each particle to radiate a small secondary wave in all directions, like a dipole antenna. All these waves add up to give specular reflection and refraction, according to the Huygens–Fresnel principle.
In the case of dielectrics such as glass, the electric field of the light acts on the electrons in the material, and the moving electrons generate fields and become new radiators. The refracted light in the glass is the combination of the forward radiation of the electrons and the incident light. The reflected light is the combination of the backward radiation of all of the electrons.
In metals, electrons with no binding energy are called free electrons. When these electrons oscillate with the incident light, the phase difference between their radiation field and the incident field is π (180°), so the forward radiation cancels the incident light, and backward radiation is just the reflected light.
Light–matter interaction in terms of photons is a topic of quantum electrodynamics, and is described in detail by Richard Feynman in his popular book QED: The Strange Theory of Light and Matter.
Diffuse reflection
When light strikes the surface of a (non-metallic) material it bounces off in all directions due to multiple reflections by the microscopic irregularities inside the material (e.g. the grain boundaries of a polycrystalline material, or the cell or fiber boundaries of an organic material) and by its surface, if it is rough. Thus, an 'image' is not formed. This is called diffuse reflection. The exact form of the reflection depends on the structure of the material. One common model for diffuse reflection is Lambertian reflectance, in which the light is reflected with equal luminance (in photometry) or radiance (in radiometry) in all directions, as defined by Lambert's cosine law.
The light sent to our eyes by most of the objects we see is due to diffuse reflection from their surface, so that this is our primary mechanism of physical observation.
Retroreflection
Some surfaces exhibit retroreflection. The structure of these surfaces is such that light is returned in the direction from which it came.
When flying over clouds illuminated by sunlight the region seen around the aircraft's shadow will appear brighter, and a similar effect may be seen from dew on grass. This partial retro-reflection is created by the refractive properties of the curved droplet's surface and reflective properties at the backside of the droplet.
Some animals' retinas act as retroreflectors (see tapetum lucidum'' for more detail), as this effectively improves the animals' night vision. Since the lenses of their eyes modify reciprocally the paths of the incoming and outgoing light the effect is that the eyes act as a strong retroreflector, sometimes seen at night when walking in wildlands with a flashlight.
A simple retroreflector can be made by placing three ordinary mirrors mutually perpendicular to one another (a corner reflector). The image produced is the inverse of one produced by a single mirror.
A surface can be made partially retroreflective by depositing a layer of tiny refractive spheres on it or by creating small pyramid like structures. In both cases internal reflection causes the light to be reflected back to where it originated. This is used to make traffic signs and automobile license plates reflect light mostly back in the direction from which it came. In this application perfect retroreflection is not desired, since the light would then be directed back into the headlights of an oncoming car rather than to the driver's eyes.
Multiple reflections
When light reflects off a mirror, one image appears. Two mirrors placed exactly face to face give the appearance of an infinite number of images along a straight line. The multiple images seen between two mirrors that sit at an angle to each other lie over a circle. The center of that circle is located at the imaginary intersection of the mirrors. A square of four mirrors placed face to face give the appearance of an infinite number of images arranged in a plane. The multiple images seen between four mirrors assembling a pyramid, in which each pair of mirrors sits an angle to each other, lie over a sphere. If the base of the pyramid is rectangle shaped, the images spread over a section of a torus.
Note that these are theoretical ideals, requiring perfect alignment of perfectly smooth, perfectly flat perfect reflectors that absorb none of the light. In practice, these situations can only be approached but not achieved because the effects of any surface imperfections in the reflectors propagate and magnify, absorption gradually extinguishes the image, and any observing equipment (biological or technological) will interfere.
Complex conjugate reflection
In this process (which is also known as phase conjugation), light bounces exactly back in the direction from which it came due to a nonlinear optical process. Not only the direction of the light is reversed, but the actual wavefronts are reversed as well. A conjugate reflector can be used to remove aberrations from a beam by reflecting it and then passing the reflection through the aberrating optics a second time. If one were to look into a complex conjugating mirror, it would be black because only the photons which left the pupil would reach the pupil.
Other types of reflection
Neutron reflection
Materials that reflect neutrons, for example beryllium, are used in nuclear reactors and nuclear weapons. In the physical and biological sciences, the reflection of neutrons off atoms within a material is commonly used to determine the material's internal structure.
Sound reflection
When a longitudinal sound wave strikes a flat surface, sound is reflected in a coherent manner provided that the dimension of the reflective surface is large compared to the wavelength of the sound. Note that audible sound has a very wide frequency range (from 20 to about 17000 Hz), and thus a very wide range of wavelengths (from about 20 mm to 17 m). As a result, the overall nature of the reflection varies according to the texture and structure of the surface. For example, porous materials will absorb some energy, and rough materials (where rough is relative to the wavelength) tend to reflect in many directions—to scatter the energy, rather than to reflect it coherently. This leads into the field of architectural acoustics, because the nature of these reflections is critical to the auditory feel of a space.
In the theory of exterior noise mitigation, reflective surface size mildly detracts from the concept of a noise barrier by reflecting some of the sound into the opposite direction. Sound reflection can affect the acoustic space.
Seismic reflection
Seismic waves produced by earthquakes or other sources (such as explosions) may be reflected by layers within the Earth. Study of the deep reflections of waves generated by earthquakes has allowed seismologists to determine the layered structure of the Earth. Shallower reflections are used in reflection seismology to study the Earth's crust generally, and in particular to prospect for petroleum and natural gas deposits.
| Physical sciences | Optics | null |
521369 | https://en.wikipedia.org/wiki/Proboscis | Proboscis | A proboscis () is an elongated appendage from the head of an animal, either a vertebrate or an invertebrate. In invertebrates, the term usually refers to tubular mouthparts used for feeding and sucking. In vertebrates, a proboscis is an elongated nose or snout.
Etymology
First attested in English in 1609 from Latin , the latinisation of the Ancient Greek (), which comes from () 'forth, forward, before' + (), 'to feed, to nourish'. The plural as derived from the Greek is , but in English the plural form proboscises occurs frequently.
Invertebrates
The most common usage is to refer to the tubular feeding and sucking organ of certain invertebrates such as insects (e.g., moths, butterflies, and mosquitoes), worms (including Acanthocephala, proboscis worms) and gastropod molluscs.
Acanthocephala
The Acanthocephala, the thorny-headed worms or spiny-headed worms, are characterized by the presence of an eversible proboscis, armed with spines, which they use to pierce and hold the gut wall of their host.
Lepidoptera mouth parts
The mouth parts of Lepidoptera (butterflies and moths) mainly consist of the sucking kind; this part is known as the proboscis or 'haustellum'. The proboscis consists of two tubes held together by hooks and separable for cleaning. The proboscis contains muscles for operating. Each tube is inwardly concave, thus forming a central tube up which moisture is sucked. Suction takes place due to the contraction and expansion of a sac in the head. A specific example of the proboscis being used for feeding is in the species Deilephila elpenor. In this species, the moth hovers in front of the flower and extends its long proboscis to attain its food.
A few Lepidoptera species lack mouth parts and therefore do not feed in the imago. Others, such as the family Micropterigidae, have mouth parts of the chewing kind.
The study of insect mouthparts was helpful for the understanding of the functional mechanism of the proboscis of butterflies (Lepidoptera) to elucidate the evolution of new form-function. The study of the proboscis of butterflies revealed surprising examples of adaptations to different kinds of fluid food, including nectar, plant sap, tree sap, dung and of adaptations to the use of pollen as complementary food in Heliconius butterflies. An extremely long proboscis appears within different groups of flower-visiting insects, but is relatively rare.
Gastropods
Some evolutionary lineages of gastropods have evolved a proboscis. In gastropods, the proboscis is an elongation of the snout with the ability to retract inside the body; it can be used for feeding, sensing the environment, and in some cases, capturing prey or attaching to hosts. Three major types of proboscises have been identified: pleurembolic (partially retractable), acrembolic (fully retractable), and intraembolic (variable in structure). Acrembolic proboscises are usually found in parasitic gastropods.
Vertebrates
The elephant's trunk and the tapir's elongated nose are called "proboscis", as is the snout of the male elephant seal.
Notable mammals with some form of proboscis are:
Aardvark
Anteater
Elephant
Elephant shrew
Hispaniolan solenodon
Echidna
Elephant seal
Leptictidium (extinct)
Moeritherium (extinct)
Numbat
Proboscis monkey
Saiga antelope
Members of the tapir family
The proboscis monkey is named for its enormous nose.
The human nose is sometimes called a proboscis, especially when large or prominent.
| Biology and health sciences | External anatomy and regions of the body | Biology |
521801 | https://en.wikipedia.org/wiki/Water%20heating | Water heating | Water heating is a heat transfer process that uses an energy source to heat water above its initial temperature. Typical domestic uses of hot water include cooking, cleaning, bathing, and space heating. In industry, hot water and water heated to steam have many uses.
Domestically, water is traditionally heated in vessels known as water heaters, kettles, cauldrons, pots, or coppers. These metal vessels that heat a batch of water do not produce a continual supply of heated water at a preset temperature. Rarely, hot water occurs naturally, usually from natural hot springs. The temperature varies with the consumption rate, becoming cooler as flow increases.
Appliances that provide a continual supply of hot water are called water heaters, hot water heaters, hot water tanks, boilers, heat exchangers, geysers (Southern Africa and the Arab world), or calorifiers. These names depend on region, and whether they heat potable or non-potable water, are in domestic or industrial use, and their energy source. In domestic installations, potable water heated for uses other than space heating is also called domestic hot water (DHW).
Fossil fuels (natural gas, liquefied petroleum gas, oil), or solid fuels are commonly used for heating water. These may be consumed directly or may produce electricity that, in turn, heats water. Electricity to heat water may also come from any other electrical source, such as nuclear power or renewable energy. Alternative energy such as solar energy, heat pumps, hot water heat recycling, and geothermal heating can also heat water, often in combination with backup systems powered by fossil fuels or electricity.
Densely populated urban areas of some countries provide district heating of hot water. This is especially the case in Scandinavia, Finland and Poland. District heating systems supply energy for water heating and space heating from combined heat and power (CHP) plants such as incinerators, central heat pumps, waste heat from industries, geothermal heating, and central solar heating. Actual heating of tap water is performed in heat exchangers at the consumers' premises. Generally the consumer has no in-building backup system as redundancy is usually significant on the district heating supply side.
Today, in the United States, domestic hot water used in homes is most commonly heated with natural gas, electric resistance, or a heat pump. Electric heat pump water heaters are significantly more efficient than electric resistance water heaters, but also more expensive to purchase. Some energy utilities offer their customers funding to help offset the higher first cost of energy efficient water heaters.
Types of water heating appliances
Hot water used for space heating may be heated by fossil fuels in a boiler, while potable water may be heated in a separate appliance. This is common practice in the US, especially when warm-air space heating is usually employed.
Storage water heaters (tank-type)
In household and commercial usage, most North American and Southern Asian water heaters are the tank type, also called storage water heaters. These consist of a cylindrical vessel or container that keeps water continuously hot and ready to use. Typical sizes for household use range from 75 to 400 L (20 to 100 US gallons). These may use electricity, natural gas, propane, heating oil, solar, or other energy sources. Natural gas heaters are most popular in the US and most European countries, since the gas is often conveniently piped throughout cities and towns and currently is the cheapest to use. In the United States, typical natural gas water heaters for households without unusual needs are with a burner rated at .
This is a popular arrangement where higher flow rates are required for limited periods. Water is heated in a pressure vessel that can withstand a hydrostatic pressure close to that of the incoming mains supply. A pressure reducing valve is sometimes employed to limit the pressure to a safe level for the vessel. In North America, these vessels are called hot water tanks, and may incorporate an electrical resistance heater, a heat pump, or a gas or oil burner that heats water directly.
Where hot-water space heating boilers are installed, domestic hot water cylinders are usually heated indirectly by primary water from the boiler, or by an electric immersion heater (often as backup to the boiler). In the UK these vessels are called indirect cylinders and direct cylinders, respectively. Additionally, if these cylinders form part of a sealed system, providing mains-pressure hot water, they are known as unvented cylinders. In the US, when connected to a boiler, they are called indirect-fired water heaters.
Compared to tankless heaters, storage water heaters have the advantage of using energy (gas or electricity) at a relatively slow rate, storing the heat for later use. The disadvantage is that over time, heat escapes through the tank wall and the water cools down, activating the heating system to heat the water back up, so investing in a tank with better insulation improves this standby efficiency. Additionally, when heavy use exhausts the hot water, there is a significant delay before hot water is available again. Larger tanks tend to provide hot water with less temperature fluctuation at moderate flow rates.
Volume storage water heaters in the United States and New Zealand are typically vertical cylindrical tanks, usually standing on the floor, a 'cylinder tray' or on a platform raised a short distance above the floor. Volume storage water heaters in Spain are typically horizontal. In India, they are mainly vertical. In apartments they can be mounted in the ceiling space over laundry-utility rooms. In Australia, gas and electric outdoor tank heaters have mainly been used (with high temperatures to increase effective capacity), but solar roof tanks are becoming fashionable.
Tiny point-of-use (POU) electric storage water heaters with capacities ranging from 832 L (26 gallons) are made for installation in kitchen and bath cabinets or on the wall above a sink. They typically use low power heating elements, about 1 kW to 1.5 kW, and can provide hot water long enough for hand washing, or, if plumbed into an existing hot water line, until hot water arrives from a remote high capacity water heater. They may be used when retrofitting a building with hot water plumbing is too costly or impractical. Since they maintain water temperature thermostatically, they can only supply a continuous flow of hot water at extremely low flow rates, unlike high-capacity tankless heaters.
In tropical countries like Singapore and India, a storage water heater may vary from 10 L to 35 L. Smaller water heaters are sufficient, as ambient weather temperatures and incoming water temperature are moderate. The Coldest regions in India like Kashmir, people are mostly dependent on the storage type electric water heaters. Mostly 50L or 75L Storage type electric water heaters are connected to overhead water source.
Point-of-use (POU) vs centralized hot water
A locational design decision may be made between point-of-use and centralized water heaters. Centralized water heaters are more traditional, and are still a good choice for small buildings. For larger buildings with intermittent or occasional hot water use, multiple POU water heaters may be a better choice, since they can reduce long waits for hot water to arrive from a remote heater. The decision where to locate the water heater(s) is only partially independent of the decision of a tanked vs. tankless water heater, or the choice of energy source for the heat.
Instantaneous water heaters (tankless-type)
Tankless water heaters—also called instantaneous, continuous flow, inline, flash, on-demand, or instant-on water heaters—are gaining in popularity. These high-power water heaters instantly heat water as it flows through the device, and do not retain any water internally except for what is in the heat exchanger coil. Copper heat exchangers are preferred in these units because of their high thermal conductivity and ease of fabrication.
Tankless heaters may be installed throughout a household at more than one point-of-use (POU), far from a central water heater, or larger centralized models may still be used to provide all the hot water requirements for an entire house. The main advantages of tankless water heaters are a plentiful continuous flow of hot water (as compared to a limited flow of continuously heated hot water from conventional tank water heaters), and potential energy savings under some conditions. The main disadvantage is their much higher initial costs; a US study in Minnesota reported a 20- to 40-year payback for the tankless water heaters. In a comparison to a less efficient natural gas fired hot water tank, on-demand natural gas will cost 30% more over its useful life.
Stand-alone appliances for quickly heating water for domestic usage are known in North America as tankless or on demand water heaters. In some places, they are called multipoint heaters, geysers or ascots. In Australia and New Zealand they are called instantaneous hot water units. In Argentina they are called calefones. In that country calefones use gas instead of electricity, although gas powered tankless water heaters can also be found in other countries. A similar wood-fired appliance was known as the chip heater.
A common arrangement where hot-water space heating is employed is for a boiler also to heat potable water, providing a continuous supply of hot water without extra equipment. Appliances that can supply both space-heating and domestic hot water are called combination (or combi) boilers. Though on-demand heaters provide a continuous supply of domestic hot water, the rate at which they can produce it is limited by the thermodynamics of heating water from the available fuel supplies.
Electric shower heads
Solar water heaters
Increasingly, solar powered water heaters are being used. Their solar collectors are installed outside dwellings, typically on the roof or walls or nearby, and the potable hot water storage tank is typically a pre-existing or new conventional water heater, or a water heater specifically designed for solar thermal. In Cyprus and Israel 90 percent of homes have solar water heating systems.
The most basic solar thermal models are the direct-gain type, in which the potable water is directly sent into the collector. Many such systems are said to use integrated collector storage (ICS), as direct-gain systems typically have storage integrated within the collector. Heating water directly is inherently more efficient than heating it indirectly via heat exchangers, but such systems offer very limited freeze protection (if any), can easily heat water to temperatures unsafe for domestic use, and ICS systems suffer from severe heat loss on cold nights and cold, cloudy days.
By contrast, indirect or closed-loop systems do not allow potable water through the panels, but rather pump a heat transfer fluid (either water or a water/antifreeze mix) through the panels. After collecting heat in the panels, the heat transfer fluid flows through a heat exchanger, transferring its heat to the potable hot water. When the panels are cooler than the storage tank or when the storage tank has already reached its maximum temperature, the controller in closed-loop systems stops the circulation pumps. In a drainback system, the water drains into a storage tank contained in conditioned or semi-conditioned space, protected from freezing temperatures. With antifreeze systems, however, the pump must be run if the panel temperature gets too hot (to prevent degradation of the antifreeze) or too cold (to prevent the water/antifreeze mixture from freezing.)
Flat panel collectors are typically used in closed-loop systems. Flat panels, which often resemble skylights, are the most durable type of collector, and they also have the best performance for systems designed for temperatures within of ambient temperature. Flat panels are regularly used in both pure water and antifreeze systems.
Another type of solar collector is the evacuated tube collector, which are intended for cold climates that do not experience severe hail and/or applications where high temperatures are needed (i.e., over ). Placed in a rack, evacuated tube collectors form a row of glass tubes, each containing absorption fins attached to a central heat-conducting rod (copper or condensation-driven). The evacuated description refers to the vacuum created in the glass tubes during the manufacturing process, which results in very low heat loss and lets evacuated tube systems achieve extreme temperatures, far in excess of water's boiling point.
Geothermal heating
In countries like Iceland and New Zealand, and other volcanic regions, water heating may be done using geothermal heating, rather than combustion.
Gravity-fed system
Where a space-heating water boiler is employed, the traditional arrangement in the UK and Ireland is to use boiler-heated (primary) water to heat potable (secondary) water contained in a cylindrical vessel (usually made of copper)—which is supplied from a cold water storage vessel or container, usually in the roof space of the building. This produces a fairly steady supply of DHW (domestic hot water) at low static pressure head but usually with a good flow. In most other parts of the world, water heating appliances do not use a cold water storage vessel or container, but heat water at pressures close to that of the incoming mains water supply.
Other improvements
Other improvements to water heaters include check valve devices at their inlet and outlet, cycle timers, electronic ignition in the case of fuel-using models, sealed air intake systems in the case of fuel-using models, and pipe insulation. The sealed air-intake system types are sometimes called "band-joist" intake units. "High-efficiency" condensing units can convert up to 98% of the energy in the fuel to heating the water. The exhaust gases of combustion are cooled and are mechanically ventilated either through the roof or through an exterior wall. At high combustion efficiencies a drain must be supplied to handle the water condensed out of the combustion products, which are primarily carbon dioxide and water vapor.
In traditional plumbing in the UK, the space-heating boiler is set up to heat a separate hot water cylinder or water heater for potable hot water. Such water heaters are often fitted with an auxiliary electrical immersion heater for use if the boiler is out of action for a time. Heat from the space-heating boiler is transferred to the water heater vessel/container by means of a heat exchanger, and the boiler operates at a higher temperature than the potable hot water supply. Most potable water heaters in North America are completely separate from the space heating units, due to the popularity of HVAC/forced air systems in North America.
Residential combustion water heaters manufactured since 2003 in the United States have been redesigned to resist ignition of flammable vapors and incorporate a thermal cutoff switch, per ANSI Z21.10.1. The first feature attempts to prevent vapors from flammable liquids and gases in the vicinity of the heater from being ignited and thus causing a house fire or explosion. The second feature prevents tank overheating due to unusual combustion conditions. These safety requirements were made in response to homeowners storing, or spilling, gasoline or other flammable liquids near their water heaters and causing fires. Since most of the new designs incorporate some type of flame arrestor screen, they require monitoring to make sure they do not become clogged with lint or dust, reducing the availability of air for combustion. If the flame arrestor becomes clogged, the thermal cutoff may act to shut down the heater.
A wetback stove (NZ), wetback heater (NZ), or back boiler (UK), is a simple household secondary water heater using incidental heat. It typically consists of a hot water pipe running behind a fireplace or stove (rather than hot water storage), and has no facility to limit the heating. Modern wetbacks may run the pipe in a more sophisticated design to assist heat-exchange. These designs are being forced out by government efficiency regulations that do not count the energy used to heat water as 'efficiently' used.
History
Another type of water heater developed in Europe predated the storage model. In London, England, in 1868, Benjamin Waddy Maughan, a painter, invented the first instantaneous domestic water heater that did not use solid fuel. Named the geyser after an Icelandic gushing hot spring, Maughan's invention made cold water at the top flow through pipes that were heated by hot gases from a burner at the bottom. Hot water then flowed into a sink or tub. The invention was somewhat dangerous because there was no flue to remove heated gases from the bathroom. A water heater is still sometimes called a geyser in the UK and South Africa.
Maughn's invention influenced the work of a Norwegian mechanical engineer named Edwin Ruud. The first automatic, storage tank-type gas water heater was invented around 1889 by Ruud after he immigrated to Pittsburgh, Pennsylvania (US). The Ruud Manufacturing Company, still in existence today, made many advancements in tank-type and tankless water heater design and operation.
Thermodynamics and economics
Water typically enters residences in the US at about , depending on latitude and season. Hot water temperatures of are usual for dish-washing, laundry and showering, which requires that the heater raise the water temperature about if the hot water is mixed with cold water at the point of use. The Uniform Plumbing Code reference shower flow rate is per minute. Sink and dishwasher usages range from per minute.
Natural gas is often measured by volume or heat content. Common units of measurement by volume are cubic metre or cubic feet at standard conditions or by heat content in kilowatt hours, British thermal units (BTU) or therm, which is equal to 100,000 BTU. A BTU is the energy required to raise one pound of water by one degree Fahrenheit. A US gallon of water weighs . To raise of water from to at 90% efficiency requires . A heater, as might exist in a tankless heater, would take about 15 minutes to do this. At $1 per therm, the cost of the gas would be about 40 cents. In comparison, a typical tank electric water heater has a heating element, which at 100% efficient results in a heating time of about 2.34 hours. At $0.16/kWh the electricity would cost $1.68.
Energy efficiencies of water heaters in residential use can vary greatly, particularly depending on manufacturer and model. However, electric heaters tend to be slightly more efficient (not counting power station losses) with recovery efficiency (how efficiently energy transfers to the water) reaching about 98%. Gas-fired heaters have maximum recovery efficiencies of only about 8294% (the remaining heat is lost with the flue gasses). Overall energy factors can be as low as 80% for electric and 50% for gas systems. Natural gas and propane tank water heaters with energy factors of 62% or greater, as well as electric tank water heaters with energy factors of 93% or greater, are considered high-efficiency units. Energy Star-qualified natural gas and propane tank water heaters (as of September 2010) have energy factors of 67% or higher, which is usually achieved using an intermittent pilot together with an automatic flue damper, baffle blowers, or power venting.
Direct electric resistance tank water heaters are not included in the Energy Star program; however, the Energy Star program does include electric heat pump units with energy factors of 200% or higher. Tankless gas water heaters (as of 2015) must have an energy factor of 90% or higher for Energy Star qualification. Since electricity production in thermal plants has efficiency levels ranging from only 15% to slightly over 55% (combined cycle gas turbine), with around 40% typical for thermal power stations, direct resistance electric water heating may be the least energy efficient option.
However, use of a heat pump can make electric water heaters much more energy efficient and lead to a decrease in carbon dioxide emissions, even more so if a low carbon source of electricity is used. Using district heating utilizing waste heat from electricity generation and other industries to heat residences and hot water gives an increased overall efficiency, removing the need for burning fossil fuel or using high energy value electricity to produce heat in the individual home.
Fundamentally, it takes a great deal of energy to heat water, as one may experience when waiting to boil a gallon of water on a stove. For this reason, tankless on-demand water heaters require a powerful energy source. A standard 120V, 15-ampere rated wall electric outlet, by comparison, only sources enough power to warm a disappointingly small amount of water: about per minute at temperature elevation.
The energy used by an electric water heater can be reduced by as much as 18% through optimal schedule and temperature control that is based on knowledge of the usage pattern.
US minimum requirements
On April 16, 2015, as part of the National Appliance Energy Conservation Act (NAECA), new minimum standards for efficiency of residential water heaters set by the United States Department of Energy went into effect. All new gas storage tank water heaters with capacities smaller than sold in the United States in 2015 or later shall have an energy factor of at least 60% (for 50-US-gallon units, higher for smaller units), increased from the pre-2015 minimum standard of 58% energy factor for 50-US-gallon gas units. Electric storage tank water heaters with capacities less than 55 US gallons sold in the United States shall have an energy factor of at least 95%, increased from the pre-2015 minimum standard of 90% for 50-US-gallon electric units.
Under the 2015 standard, for the first time, storage water heaters with capacities of 55 US gallons or larger now face stricter efficiency requirements than those of 50 US gallons or less. Under the pre-2015 standard, a gas storage water heater with a nominal input of or less was able to have an energy factor as low as 53%, while under the 2015 standard, the minimum energy factor for a 75-US-gallon gas storage tank water heater is now 74%, which can only be achieved by using condensing technology. Storage water heaters with a nominal input of or greater are not currently affected by these requirements, since energy factor is not defined for such units. An electric storage tank water heater was able to have a minimum energy factor of 86% under the pre-2015 standard, while under the 2015 standard, the minimum energy factor for an 80-gallon electric storage tank water heater is now 197%, which is only possible with heat pump technology. This rating measures efficiency at the point of use.
Depending on how electricity is generated, overall efficiency may be much lower. For example, in a traditional coal plant, only about 30–35% of the energy in the coal ends up as electricity on the other end of the generator. Losses on the electrical grid (including line losses and voltage transformation losses) reduce electrical efficiency further. According to data from the Energy Information Administration, transmission and distribution losses in 2005 consumed 6.1% of net generation. In contrast, 90% of natural gas's energy value is
delivered to the consumer. (In neither case is the energy expended exploring, developing and extracting coal or natural gas resources included in the quoted efficiency numbers.) Gas tankless water heaters shall have an energy factor of 82% or greater under the 2015 standards, which corresponds to the pre-2015 Energy Star standard.
In 2022 the Department of Energy proposed rules that would take effect in 2026 and would effectively eliminate inefficient non-condensing gas water heaters in commercial buildings. Non-condensing models waste heat, while condensing models capture and used otherwise lost energy. The change will reduce emissions by 38 million tons of carbon dioxide over 30 years and reduce buildings' energy costs.
Water heater safety
Explosion hazard
Water heaters potentially can explode and cause significant damage, injury, or death if certain safety devices are not installed. A safety device called a temperature and pressure relief (T&P or TPR) valve, is normally fitted on the top of the water heater to dump water if the temperature or pressure becomes too high. Most plumbing codes require that a discharge pipe be connected to the valve to direct the flow of discharged hot water to a drain, typically a nearby floor drain, or outside the living space. Some building codes allow the discharge pipe to terminate in the garage.
If a gas or propane fired water heater is installed in a garage or basement, many plumbing codes require that it be elevated at least above the floor to reduce the potential for fire or explosion due to spillage or leakage of combustible liquids in the garage. Furthermore, certain local codes mandate that tank-type heaters in new and retrofit installations must be secured to an adjacent wall by a strap or anchor to prevent tipping over and breaking the water and gas pipes in the event of an earthquake.
For older houses where the water heater is part of the space heating boiler, and plumbing codes allow, some plumbers install an automatic gas shutoff (such as the "Watts 210") in addition to a TPR valve. When the device senses that the temperature reaches , it shuts off the gas supply and prevents further heating. In addition, an expansion tank or exterior pressure relief valve must be installed to prevent pressure buildup in the plumbing from rupturing pipes, valves, or the water heater.
Thermal burns (scalding)
Scalding is a serious concern with any water heater. Human skin burns quickly at high temperature, in less than 5 seconds at , but much slower at — it takes a full minute for a second degree burn. Older people and children often receive serious scalds due to disabilities or slow reaction times. In the United States and elsewhere it is common practice to put a tempering valve or thermostatic mixing valve on the outlet of the water heater. The result of automatically mixing hot and cold water via a tempering valve is referred to as "tempered water".
A tempering valve mixes enough cold water with the hot water from the heater to keep the outgoing water temperature fixed at a more moderate temperature, often set to . Without a tempering valve, reduction of the water heater's setpoint temperature is the most direct way to reduce scalding. However, for sanitation, hot water is needed at a temperature that can cause scalding. This may be accomplished by using a supplemental heater in an appliance that requires hotter water.
Most residential dishwashing machines, for example, include an internal electric heating element for increasing the water temperature above that provided by a domestic water heater.
Bacterial contamination
Two conflicting safety issues affect water heater temperature—the risk of scalding from excessively hot water greater than , and the risk of incubating bacteria colonies, particularly Legionella, in water that is not hot enough to kill them. Both risks are potentially life-threatening and are balanced by setting the water heater's thermostat to . The European Guidelines for Control and Prevention of Travel Associated Legionnaires' Disease recommend that hot water should be stored at and distributed so that a temperature of at least and preferably is achieved within one minute at points of use.
If there is a dishwasher without a booster heater, it may require a water temperature within a range of for optimum cleaning, but tempering valves set to no more than can be applied to faucets to avoid scalding. Tank temperatures above may produce limescale deposits, which could later harbor bacteria, in the water tank. Higher temperatures may also increase etching of glassware in the dishwasher.
Tank thermostats are not a reliable guide to the internal temperature of the tank. Gas-fired water tanks may have no temperature calibration shown. An electric thermostat shows the temperature at the elevation of the thermostat, but water lower in the tank can be considerably cooler. An outlet thermometer is a better indication of water temperature.
In the renewable energy industry (solar and heat pumps, in particular) the conflict between daily thermal Legionella control and high temperatures, which may drop system performance, is subject to heated debate. In a paper seeking a green exemption from normal Legionellosis safety standards, Europe's top CEN solar thermal technical committee TC 312 asserts that a 50% fall in performance would occur if solar water heating systems were heated to the base daily. However some solar simulator analysis work using Polysun 5 suggests that an 11% energy penalty is a more likely figure. Whatever the context, both energy efficiency and scalding safety requirements push in the direction of considerably lower water temperatures than the legionella pasteurization temperature of around .
Legionella pneumophila has been detected at the point of use downstream from horizontally-mounted electric water heaters with volumes of 150 Liters.
However, legionella can be safely and easily controlled with good design and engineering protocols. For instance raising the temperature of water heaters once a day or even once every few days to at the coldest part of the water heater for 30 minutes effectively controls legionella. In all cases and in particular energy efficient applications, Legionnaires' disease is more often than not the result of engineering design issues that do not take into consideration the impact of stratification or low flow.
It is also possible to control Legionella risks by chemical treatment of the water. This technique allows lower water temperatures to be maintained in the pipework without the associated Legionella risk. The benefit of lower pipe temperatures is that the heat loss rate is reduced and thus the energy consumption is reduced.
| Technology | Household appliances | null |
521830 | https://en.wikipedia.org/wiki/Orthorhombic%20crystal%20system | Orthorhombic crystal system | In crystallography, the orthorhombic crystal system is one of the 7 crystal systems. Orthorhombic lattices result from stretching a cubic lattice along two of its orthogonal pairs by two different factors, resulting in a rectangular prism with a rectangular base (a by b) and height (c), such that a, b, and c are distinct. All three bases intersect at 90° angles, so the three lattice vectors remain mutually orthogonal.
Bravais lattices
There are four orthorhombic Bravais lattices: primitive orthorhombic, base-centered orthorhombic, body-centered orthorhombic, and face-centered orthorhombic.
For the base-centered orthorhombic lattice, the primitive cell has the shape of a right rhombic prism; it can be constructed because the two-dimensional centered rectangular base layer can also be described with primitive rhombic axes. Note that the length of the primitive cell below equals of the conventional cell above.
Crystal classes
The orthorhombic crystal system class names, examples, Schönflies notation, Hermann-Mauguin notation, point groups, International Tables for Crystallography space group number, orbifold notation, type, and space groups are listed in the table below.
In two dimensions
In two dimensions there are two orthorhombic Bravais lattices: primitive rectangular and centered rectangular.
| Physical sciences | Crystallography | Physics |
521843 | https://en.wikipedia.org/wiki/Hydraulic%20engineering | Hydraulic engineering | Hydraulic engineering as a sub-discipline of civil engineering is concerned with the flow and conveyance of fluids, principally water and sewage. One feature of these systems is the extensive use of gravity as the motive force to cause the movement of the fluids. This area of civil engineering is intimately related to the design of bridges, dams, channels, canals, and levees, and to both sanitary and environmental engineering.
Hydraulic engineering is the application of the principles of fluid mechanics to problems dealing with the collection, storage, control, transport, regulation, measurement, and use of water. Before beginning a hydraulic engineering project, one must figure out how much water is involved. The hydraulic engineer is concerned with the transport of sediment by the river, the interaction of the water with its alluvial boundary, and the occurrence of scour and deposition. "The hydraulic engineer actually develops conceptual designs for the various features which interact with water such as spillways and outlet works for dams, culverts for highways, canals and related structures for irrigation projects, and cooling-water facilities for thermal power plants."
Fundamental principles
A few examples of the fundamental principles of hydraulic engineering include fluid mechanics, fluid flow, behavior of real fluids, hydrology, pipelines, open channel hydraulics, mechanics of sediment transport, physical modeling, hydraulic machines, and drainage hydraulics.
Fluid mechanics
Fundamentals of Hydraulic Engineering defines hydrostatics as the study of fluids at rest. In a fluid at rest, there exists a force, known as pressure, that acts upon the fluid's surroundings. This pressure, measured in N/m2, is not constant throughout the body of fluid. Pressure, p, in a given body of fluid, increases with an increase in depth. Where the upward force on a body acts on the base and can be found by the equation:
where,
ρ = density of water
g = specific gravity
y = depth of the body of liquid
Rearranging this equation gives you the pressure head . Four basic devices for pressure measurement are a piezometer, manometer, differential manometer, Bourdon gauge, as well as an inclined manometer.
As Prasuhn states:
On undisturbed submerged bodies, pressure acts along all surfaces of a body in a liquid, causing equal perpendicular forces in the body to act against the pressure of the liquid. This reaction is known as equilibrium. More advanced applications of pressure are that on plane surfaces, curved surfaces, dams, and quadrant gates, just to name a few.
Behavior of real fluids
Real and Ideal fluids
The main difference between an ideal fluid and a real fluid is that for ideal flow p1 = p2 and for real flow p1 > p2. Ideal fluid is incompressible and has no viscosity. Real fluid has viscosity. Ideal fluid is only an imaginary fluid as all fluids that exist have some viscosity.
Viscous flow
A viscous fluid will deform continuously under a shear force by the pascles law, whereas an ideal fluid does not deform.
Laminar flow and turbulence
The various effects of disturbance on a viscous flow are a stable, transition and unstable.
Bernoulli's equation
For an ideal fluid, Bernoulli's equation holds along streamlines.
As the flow comes into contact with the plate, the layer of fluid actually "adheres" to a solid surface. There is then a considerable shearing action between the layer of fluid on the plate surface and the second layer of fluid. The second layer is therefore forced to decelerate (though it is not quite brought to rest), creating a shearing action with the third layer of fluid, and so on. As the fluid passes further along with the plate, the zone in which shearing action occurs tends to spread further outwards. This zone is known as the "boundary layer". The flow outside the boundary layer is free of shear and viscous-related forces so it is assumed to act as an ideal fluid. The intermolecular cohesive forces in a fluid are not great enough to hold fluid together. Hence a fluid will flow under the action of the slightest stress and flow will continue as long as the stress is present. The flow inside the layer can be either vicious or turbulent, depending on Reynolds number.
Applications
Common topics of design for hydraulic engineers include hydraulic structures such as dams, levees, water distribution networks including both domestic and fire water supply, distribution and automatic sprinkler systems, water collection networks, sewage collection networks, storm water management, sediment transport, and various other topics related to transportation engineering and geotechnical engineering. Equations developed from the principles of fluid dynamics and fluid mechanics are widely utilized by other engineering disciplines such as mechanical, aeronautical and even traffic engineers.
Related branches include hydrology and rheology while related applications include hydraulic modeling, flood mapping, catchment flood management plans, shoreline management plans, estuarine strategies, coastal protection, and flood alleviation.
History
Antiquity
Earliest uses of hydraulic engineering were to irrigate crops and dates back to the Middle East and Africa. Controlling the movement and supply of water for growing food has been used for many thousands of years. One of the earliest hydraulic machines, the water clock was used in the early 2nd millennium BC. Other early examples of using gravity to move water include the Qanat system in ancient Persia and the very similar Turpan water system in ancient China as well as irrigation canals in Peru.
In ancient China, hydraulic engineering was highly developed, and engineers constructed massive canals with levees and dams to channel the flow of water for irrigation, as well as locks to allow ships to pass through. Sunshu Ao is considered the first Chinese hydraulic engineer. Another important Hydraulic Engineer in China, Ximen Bao was credited of starting the practice of large scale canal irrigation during the Warring States period (481 BC–221 BC), even today hydraulic engineers remain a respectable position in China.
In the Archaic epoch of the Philippines, hydraulic engineering also developed specially in the Island of Luzon, the Ifugaos of the mountainous region of the Cordilleras built irrigations, dams and hydraulic works and the famous Banaue Rice Terraces as a way for assisting in growing crops around 1000 BC. These Rice Terraces are 2,000-year-old terraces that were carved into the mountains of Ifugao in the Philippines by ancestors of the indigenous people. The Rice Terraces are commonly referred to as the "Eighth Wonder of the World". It is commonly thought that the terraces were built with minimal equipment, largely by hand. The terraces are located approximately 1500 metres (5000 ft) above sea level. They are fed by an ancient irrigation system from the rainforests above the terraces. It is said that if the steps were put end to end, it would encircle half the globe.
Eupalinos of Megara was an ancient Greek engineer who built the Tunnel of Eupalinos on Samos in the 6th century BC, an important feat of both civil and hydraulic engineering. The civil engineering aspect of this tunnel was that it was dug from both ends which required the diggers to maintain an accurate path so that the two tunnels met and that the entire effort maintained a sufficient slope to allow the water to flow.
Hydraulic engineering was highly developed in Europe under the aegis of the Roman Empire where it was especially applied to the construction and maintenance of aqueducts to supply water to and remove sewage from their cities. In addition to supplying the needs of their citizens they used hydraulic mining methods to prospect and extract alluvial gold deposits in a technique known as hushing, and applied the methods to other ores such as those of tin and lead.
In the 15th century, the Somali Ajuran Empire was the only hydraulic empire in Africa. As a hydraulic empire, the Ajuran State monopolized the water resources of the Jubba and Shebelle Rivers. Through hydraulic engineering, it also constructed many of the limestone wells and cisterns of the state that are still operative and in use today. The rulers developed new systems for agriculture and taxation, which continued to be used in parts of the Horn of Africa as late as the 19th century.
Further advances in hydraulic engineering occurred in the Muslim world between the 8th and 16th centuries, during what is known as the Islamic Golden Age. Of particular importance was the 'water management technological complex' which was central to the Islamic Green Revolution. The various components of this 'toolkit' were developed in different parts of the Afro-Eurasian landmass, both within and beyond the Islamic world. However, it was in the medieval Islamic lands where the technological complex was assembled and standardized, and subsequently diffused to the rest of the Old World. Under the rule of a single Islamic caliphate, different regional hydraulic technologies were assembled into "an identifiable water management technological complex that was to have a global impact." The various components of this complex included canals, dams, the qanat system from Persia, regional water-lifting devices such as the noria, shaduf and screwpump from Egypt, and the windmill from Islamic Afghanistan. Other original Islamic developments included the saqiya with a flywheel effect from Islamic Spain, the reciprocating suction pump and crankshaft-connecting rod mechanism from Iraq, and the geared and hydropowered water supply system from Syria.
Modern times
In many respects, the fundamentals of hydraulic engineering have not changed since ancient times. Liquids are still moved for the most part by gravity through systems of canals and aqueducts, though the supply reservoirs may now be filled using pumps. The need for water has steadily increased from ancient times and the role of the hydraulic engineer is a critical one in supplying it. For example, without the efforts of people like William Mulholland the Los Angeles area would not have been able to grow as it has because it simply does not have enough local water to support its population. The same is true for many of our world's largest cities. In much the same way, the central valley of California could not have become such an important agricultural region without effective water management and distribution for irrigation. In a somewhat parallel way to what happened in California, the creation of the Tennessee Valley Authority (TVA) brought work and prosperity to the South by building dams to generate cheap electricity and control flooding in the region, making rivers navigable and generally modernizing life in the region.
Leonardo da Vinci (1452–1519) performed experiments, investigated and speculated on waves and jets, eddies and streamlining. Isaac Newton (1642–1727) by formulating the laws of motion and his law of viscosity, in addition to developing the calculus, paved the way for many great developments in fluid mechanics. Using Newton's laws of motion, numerous 18th-century mathematicians solved many frictionless (zero-viscosity) flow problems. However, most flows are dominated by viscous effects, so engineers of the 17th and 18th centuries found the inviscid flow solutions unsuitable, and by experimentation they developed empirical equations, thus establishing the science of hydraulics.
Late in the 19th century, the importance of dimensionless numbers and their relationship to turbulence was recognized, and dimensional analysis was born. In 1904 Ludwig Prandtl published a key paper, proposing that the flow fields of low-viscosity fluids be divided into two zones, namely a thin, viscosity-dominated boundary layer near solid surfaces, and an effectively inviscid outer zone away from the boundaries. This concept explained many former paradoxes and enabled subsequent engineers to analyze far more complex flows. However, we still have no complete theory for the nature of turbulence, and so modern fluid mechanics continues to be combination of experimental results and theory.
The modern hydraulic engineer uses the same kinds of computer-aided design (CAD) tools as many of the other engineering disciplines while also making use of technologies like computational fluid dynamics to perform the calculations to accurately predict flow characteristics, GPS mapping to assist in locating the best paths for installing a system and laser-based surveying tools to aid in the actual construction of a system.
| Technology | Disciplines | null |
521877 | https://en.wikipedia.org/wiki/Sodium%20hydride | Sodium hydride | Sodium hydride is the chemical compound with the empirical formula NaH. This alkali metal hydride is primarily used as a strong yet combustible base in organic synthesis. NaH is a saline (salt-like) hydride, composed of Na+ and H− ions, in contrast to molecular hydrides such as borane, silane, germane, ammonia, and methane. It is an ionic material that is insoluble in all solvents (other than molten sodium metal), consistent with the fact that H− ions do not exist in solution.
Basic properties and structure
NaH is colorless, although samples generally appear grey. NaH is around 40% denser than Na (0.968 g/cm3).
NaH, like LiH, KH, RbH, and CsH, adopts the NaCl crystal structure. In this motif, each Na+ ion is surrounded by six H− centers in an octahedral geometry. The ionic radii of H− (146 pm in NaH) and F− (133 pm) are comparable, as judged by the Na−H and Na−F distances.
"Inverse sodium hydride" (hydrogen sodide)
A very unusual situation occurs in a compound dubbed "inverse sodium hydride", which contains H+ and Na− ions. Na− is an alkalide, and this compound differs from ordinary sodium hydride in having a much higher energy content due to the net displacement of two electrons from hydrogen to sodium. A derivative of this "inverse sodium hydride" arises in the presence of the base [36]adamanzane. This molecule irreversibly encapsulates the H+ and shields it from interaction with the alkalide Na−. Theoretical work has suggested that even an unprotected protonated tertiary amine complexed with the sodium alkalide might be metastable under certain solvent conditions, though the barrier to reaction would be small and finding a suitable solvent might be difficult.
Preparation
Industrially, NaH is prepared by introducing molten sodium into mineral oil with hydrogen at atmospheric pressure and mixed vigorously at ~8000 rpm. The reaction is especially rapid at 250−300 °C.
The resultant suspension of NaH in mineral oil is often directly used, such as in the production of diborane.
Applications in organic synthesis
As a strong base
NaH is a base of wide scope and utility in organic chemistry. As a superbase, it is capable of deprotonating a range of even weak Brønsted acids to give the corresponding sodium derivatives. Typical "easy" substrates contain O-H, N-H, S-H bonds, including alcohols, phenols, pyrazoles, and thiols.
NaH notably deprotonates carbon acids (i.e., C-H bonds) such as 1,3-dicarbonyls such as malonic esters. The resulting sodium derivatives can be alkylated. NaH is widely used to promote condensation reactions of carbonyl compounds via the Dieckmann condensation, Stobbe condensation, Darzens condensation, and Claisen condensation. Other carbon acids susceptible to deprotonation by NaH include sulfonium salts and DMSO. NaH is used to make sulfur ylides, which in turn are used to convert ketones into epoxides, as in the Johnson–Corey–Chaykovsky reaction.
As a reducing agent
NaH reduces certain main group compounds, but analogous reactivity is very rare in organic chemistry (see below). Notably boron trifluoride reacts to give diborane and sodium fluoride:
6 NaH + 2 BF3 → B2H6 + 6 NaF
Si–Si and S–S bonds in disilanes and disulfides are also reduced.
A series of reduction reactions, including the hydrodecyanation of tertiary nitriles, reduction of imines to amines, and amides to aldehydes, can be effected by a composite reagent composed of sodium hydride and an alkali metal iodide (NaH⋅MI, M = Li, Na).
Hydrogen storage
Although not commercially significant sodium hydride has been proposed for hydrogen storage for use in fuel cell vehicles. In one experimental implementation, plastic pellets containing NaH are crushed in the presence of water to release the hydrogen. One challenge with this technology is the regeneration of NaH from the NaOH formed by hydrolysis.
Practical considerations
Sodium hydride is sold as a mixture of 60% sodium hydride (w/w) in mineral oil. Such a dispersion is safer to handle and weigh than pure NaH. The compound is often used in this form but the pure grey solid can be prepared by rinsing the commercial product with pentane or tetrahydrofuran, with care being taken because the waste solvent will contain traces of NaH and can ignite in air. Reactions involving NaH usually require air-free techniques.
Safety
NaH can ignite spontaneously in air. It also reacts vigorously with water or humid air to release hydrogen, which is very flammable, and sodium hydroxide (NaOH), a quite corrosive base. In practice, most sodium hydride is sold as a dispersion in mineral oil, which can be safely handled in air. Although sodium hydride is widely used in DMSO, DMF or DMAc for SN2 type reactions there have been many cases of fires and/or explosions from such mixtures.
| Physical sciences | Hydride salts | Chemistry |
521948 | https://en.wikipedia.org/wiki/Zingiber | Zingiber | Zingiber is a genus of flowering plants in the family Zingiberaceae. It is native to China, the Indian Subcontinent, New Guinea, and Southeast Asia, especially Thailand. It contains the true gingers, plants grown the world over for their culinary value. The most well known species are Z. officinale and Z. mioga, two garden gingers.
Culinary
Each ginger species has a different culinary usage; for example, myoga is valued for the stem and flowers. Garden ginger's rhizome is the classic spice "ginger", and may be used whole, candied (known commonly as crystallized ginger), or dried and powdered. Other popular gingers used in cooking include cardamom and turmeric, though neither of these examples is a "true ginger" – they belong to different genera in the family Zingiberaceae.
Species
Plants of the World Online currently includes:
Zingiber acuminatum Valeton
Zingiber aguingayae Docot
Zingiber albiflorum R.M.Sm.
Zingiber album Nurainas
Zingiber anamalayanum Sujanapal & Sasidh.
Zingiber angustifolium C.K.Lim & Meekiong
Zingiber apoense Elmer
Zingiber argenteum Mood & Theilade
Zingiber arunachalensis A.Joe, T.Jayakr., Hareesh & M.Sabu
Zingiber atroporphyreum Škornick. & Q.B.Nguyen
Zingiber atrorubens Gagnep.
Zingiber aurantiacum (Holttum) Theilade
Zingiber banhaoense Mood & Theilade
Zingiber barbatum Wall.
Zingiber belumense C.K.Lim & Meekiong
Zingiber bipinianum D.K.Roy, D.Verma, Talukdar & Dutta Choud.
Zingiber bisectum D.Fang
Zingiber brachystachys Triboun & K.Larsen
Zingiber bradleyanum Craib
Zingiber brevifolium N.E.Br.
Zingiber bulusanense Elmer
Zingiber callianthus Triboun & K.Larsen
Zingiber capitatum Roxb.
Zingiber cardiocheilum Škornick. & Q.B.Nguyen
Zingiber castaneum Škornick. & Q.B.Nguyen
Zingiber caudatum Biseshwori & Bipin
Zingiber cernuum Dalzell
Zingiber chantaranothaii Triboun & K.Larsen
Zingiber chengii Y.H.Tseng, C.M.Wang & Y.C.Lin
Zingiber chlorobracteatum Mood & Theilade
Zingiber chrysanthum Roscoe
Zingiber chrysostachys Ridl.
Zingiber citriodorum Theilade & Mood
Zingiber clarkei King ex Baker
Zingiber cochleariforme D.Fang
Zingiber collinsii Mood & Theilade
Zingiber coloratum N.E.Br.
Zingiber corallinum Hance
Zingiber cornubracteatum Triboun & K.Larsen
Zingiber curtisii Holttum
Zingiber cylindricum Thwaites
Zingiber densissimum S.Q.Tong & Y.M.Xia
Zingiber discolor Škornick., H.Ð.Tran & Rybková
Zingiber diwakarianum R.Kr.Singh
Zingiber eberhardtii Gagnep.
Zingiber eborinum Mood & Theilade
Zingiber elatius (Ridl.) Theilade
Zingiber elatum Roxb.
Zingiber ellipticum (S.Q.Tong & Y.M.Xia) Q.G.Wu & T.L.Wu
Zingiber engganoense Ardiyani
Zingiber fallax (Loes.) L.Bai, Juan Chen & N.H.Xia
Zingiber flagelliforme Mood & Theilade
Zingiber flammeum Theilade & Mood
Zingiber flaviflorum C.K.Lim & Meekiong
Zingiber flavofusiforme M.M.Aung & Nob.Tanaka
Zingiber flavomaculosum S.Q.Tong
Zingiber flavovirens Theilade
Zingiber fragile S.Q.Tong
Zingiber fraseri Theilade
Zingiber georgeae Mood & Theilade
Zingiber gracile Jack
Zingiber gramineum Noronha ex Blume
Zingiber griffithii Baker
Zingiber guangxiense D.Fang
Zingiber gulinense Y.M.Xia
Zingiber hainanense Y.S.Ye, L.Bai & N.H.Xia
Zingiber idae Triboun & K.Larsen
Zingiber incomptum B.L.Burtt & R.M.Sm.
Zingiber inflexum Blume
Zingiber integrilabrum Hance
Zingiber integrum S.Q.Tong
Zingiber intermedium Baker
Zingiber isanense Triboun & K.Larsen
Zingiber jiewhoei Škornick.
Zingiber junceum Gagnep.
Zingiber kangleipakense Kishor & Škornick.
Zingiber kawagoii Hayata
Zingiber kelabitianum Theilade & H.Chr.
Zingiber kerrii Craib
Zingiber kunstleri King ex Ridl.
Zingiber lambii Mood & Theilade
Zingiber laoticum Gagnep.
Zingiber larsenii Theilade
Zingiber latifolium Theilade & Mood
Zingiber lecongkietii Škornick. & H.Ð.Tran
Zingiber leptorrhizum D.Fang
Zingiber leptostachyum Valeton
Zingiber leucochilum L.Bai, Skornick. & N.H.Xia
Zingiber ligulatum Roxb.
Zingiber limianum Meekiong
Zingiber lingyunense D.Fang
Zingiber loerzingii Valeton
Zingiber longibracteatum Theilade
Zingiber longiglande D.Fang & D.H.Qin
Zingiber longiligulatum S.Q.Tong
Zingiber longipedunculatum Ridl.
Zingiber longyanjiang Z.Y.Zhu
Zingiber macradenium K.Schum.
Zingiber macrocephalum (Zoll.) K.Schum.
Zingiber macroglossum Valeton
Zingiber macrorrhynchus K.Schum.
Zingiber malaysianum C.K.Lim : black ginger
Zingiber marginatum Roxb.
Zingiber martini R.M.Sm.
Zingiber matangense Noor Ain, Tawan & Meekiong
Zingiber matupiense M.M.Aung & Nob.Tanaka
Zingiber matutumense Mood & Theilade
Zingiber mawangense Noor Ain & Meekiong
Zingiber meghalayense Sushil K.Singh, Ram.Kumar & Mood
Zingiber mekongense Gagnep.
Zingiber mellis Škornick., H.Ð.Tran & Sída f.
Zingiber microcheilum Škornick., H.Ð.Tran & Sída f.
Zingiber mioga (Thunb.) Roscoe
Zingiber mizoramense Ram.Kumar, Sushil K.Singh & S.Sharma
Zingiber molle Ridl.
Zingiber monglaense S.J.Chen & Z.Y.Chen
Zingiber monophyllum Gagnep.
Zingiber montanum (J.Koenig) Link ex A.Dietr. (synonyms: Z. cassumunar; Z. purpureum)
Zingiber multibracteatum Holttum
Zingiber murlenica Ram.Kumar, Sushil K.Singh & S.Sharma
Zingiber nanlingense Lin Chen, A.Q.Dong & F.W.Xing
Zingiber natmataungense S.S.Zhou & R.Li
Zingiber nazrinii C.K.Lim & Meekiong
Zingiber neesanum (J.Graham) Ramamoorthy
Zingiber neglectum Valeton
Zingiber negrosense Elmer
Zingiber neotruncatum T.L.Wu, K.Larsen & Turland
Zingiber nigrimaculatum S.Q.Tong
Zingiber nimmonii (J.Graham) Dalzell
Zingiber nitens M.F.Newman
Zingiber niveum Mood & Theilade
Zingiber odoriferum Blume
Zingiber officinale Roscoe
Zingiber oligophyllum K.Schum.
Zingiber olivaceum Mood & Theilade
Zingiber orbiculatum S.Q.Tong
Zingiber ottensii Valeton
Zingiber pachysiphon B.L.Burtt & R.M.Sm.
Zingiber panduratum Roxb.
Zingiber papuanum Valeton
Zingiber pardocheilum Wall. ex Baker
Zingiber parishii Hook.f.
Zingiber pauciflorum L.Bai, Skornick., D.Z.Li & N.H.Xia
Zingiber pellitum Gagnep.
Zingiber pendulum Mood & Theilade
Zingiber petiolatum (Holttum) Theilade
Zingiber pherimaense Biseshwori & Bipin
Zingiber phillippsiae Mood & Theilade
Zingiber phumiangense Chaveer. & Mokkamul
Zingiber pleiostachyum K.Schum.
Zingiber plicatum Škornick. & Q.B.Nguyen
Zingiber popaense Nob.Tanaka
Zingiber porphyrochilum Y.H.Tan & H.B.Ding
Zingiber porphyrosphaerum K.Schum.
Zingiber pseudopungens R.M.Sm.
Zingiber pseudosquarrosum L.J.Singh & P.Singh
Zingiber puberulum Ridl.
Zingiber purpureum Roscoe
Zingiber pyroglossum Triboun & K.Larsen
Zingiber raja C.K.Lim & Kharuk.
Zingiber recurvatum S.Q.Tong & Y.M.Xia
Zingiber roseum (Roxb.) Roscoe
Zingiber rubens Roxb.
Zingiber rufopilosum Gagnep.
Zingiber sabuanum K.M.P.Kumar & A.Joe
Zingiber sabun C.K.Lim
Zingiber sadakornii Triboun & K.Larsen
Zingiber salarkhanii M.A.Rahman & Yusuf
Zingiber shuanglongense C.L.Yeh & S.W.Chung
Zingiber simaoense Y.Y.Qian
Zingiber singapurense Škornick.
Zingiber skornickovae N.S.Lý
Zingiber smilesianum Craib
Zingiber spectabile Griff.
Zingiber squarrosum Roxb.
Zingiber stenostachys K.Schum.
Zingiber striolatum Diels
Zingiber subroseum Docot
Zingiber sulphureum Burkill ex Theilade
Zingiber tenuifolium L.Bai, Škornick. & N.H.Xia
Zingiber tenuiscapus Triboun & K.Larsen
Zingiber thorelii Gagnep.
Zingiber tuanjuum Z.Y.Zhu
Zingiber ultralimitale Ardiyani & A.D.Poulsen
Zingiber vanlithianum Koord.
Zingiber velutinum Mood & Theilade
Zingiber ventricosum L.Bai, Škornick., N.H.Xia & Y.S.Ye
Zingiber vinosum Mood & Theilade
Zingiber viridiflavum Mood & Theilade
Zingiber vittacheilum Triboun & K.Larsen
Zingiber vuquangense N.S.Lý, T.H.Lê, T.H.Trinh, V.H.Nguyen & N.D.Do
Zingiber wandingense S.Q.Tong
Zingiber wightianum Thwaites
Zingiber wrayi Prain ex Ridl.
Zingiber yersinii Škornick., H.Ð.Tran & Rybková
Zingiber yingjiangense S.Q.Tong
Zingiber yunnanense S.Q.Tong & X.Z.Liu
Zingiber zerumbet (L.) Roscoe ex Sm. : shampoo ginger
| Biology and health sciences | Zingiberales | Plants |
521970 | https://en.wikipedia.org/wiki/Oxonium%20ion | Oxonium ion | In chemistry, an oxonium ion is any cation containing an oxygen atom that has three bonds and 1+ formal charge. The simplest oxonium ion is the hydronium ion ().
Alkyloxonium
Hydronium is one of a series of oxonium ions with the formula RnH3−nO+.
Oxygen is usually pyramidal with an sp3 hybridization. Those with n = 1 are called primary oxonium ions, an example being protonated alcohol (e.g. methanol). In acidic media, the oxonium functional group produced by protonating an alcohol can be a leaving group in the E2 elimination reaction. The product is an alkene. Extreme acidity, heat, and dehydrating conditions are usually required. Other hydrocarbon oxonium ions are formed by protonation or alkylation of alcohols or ethers (R−C−−R1R2).
Secondary oxonium ions have the formula R2OH+, an example being protonated ethers.
Tertiary oxonium ions have the formula R3O+, an example being trimethyloxonium. Tertiary alkyloxonium salts are useful alkylating agents. For example, triethyloxonium tetrafluoroborate ()(), a white crystalline solid, can be used, for example, to produce ethyl esters when the conditions of traditional Fischer esterification are unsuitable. It is also used for preparation of enol ethers and related functional groups.
Oxatriquinane and oxatriquinacene are unusually stable oxonium ions, first described in 2008. Oxatriquinane does not react with boiling water or with alcohols, thiols, halide ions, or amines, although it does react with stronger nucleophiles such as hydroxide, cyanide, and azide.
Oxocarbenium ions
Another class of oxonium ions encountered in organic chemistry is the oxocarbenium ions, obtained by protonation or alkylation of a carbonyl group e.g. R−C=−R′ which forms a resonance structure with the fully-fledged carbocation R−−O−R′ and is therefore especially stable:
Gold-stabilized species
An unusually stable oxonium species is the gold complex tris[triphenylphosphinegold(I)]oxonium tetrafluoroborate, [(Ph3PAu)3O][BF4], where the intramolecular aurophilic interactions between the gold atoms are believed responsible for the stabilisation of the cation. This complex is prepared by treatment of Ph3PAuCl with Ag2O in the presence of NaBF4:
3 Ph3PAuCl + Ag2O + NaBF4 → [(Ph3PAu)3O]+[BF4]− + 2 AgCl + NaCl
It has been used as a catalyst for the propargyl Claisen rearrangement.
Relevance to natural product chemistry
Complex bicyclic and tricyclic oxonium ions have been proposed as key intermediates in the biosynthesis of a series of natural products by the red algae of the genus Laurencia.
Several members of these elusive species have been prepared explicitly by total synthesis, demonstrating the possibility of their existence. The key to their successful generation was the use of a weakly coordinating anion (Krossing's anion, [Al(pftb)4]−, pftb = perfluoro-tert-butoxy) as the counteranion. As shown in the example below, this was executed by a transannular halide abstraction strategy through the reaction of the oxonium ion precursor (an organic halide) with the silver salt of the Krossing's anion Ag[Al(pftb)4]•CH2Cl2, generating the desired oxonium ion with simultaneous precipitation of inorganic silver halides. The resulting oxonium ions were characterized comprehensively by nuclear magnetic resonance spectroscopy at low temperature (−78 °C) with support from density functional theory computation.
These oxonium ions were also demonstrated to directly give rise to multiple related natural products by reacting with various nucleophiles, such as water, bromide, chloride, and acetate.
| Physical sciences | Atomic physics | Physics |
522130 | https://en.wikipedia.org/wiki/Protonation | Protonation | In chemistry, protonation (or hydronation) is the adding of a proton (or hydron, or hydrogen cation), usually denoted by H+, to an atom, molecule, or ion, forming a conjugate acid. (The complementary process, when a proton is removed from a Brønsted–Lowry acid, is deprotonation.) Some examples include
The protonation of water by sulfuric acid:
H2SO4 + H2O H3O+ +
The protonation of isobutene in the formation of a carbocation:
(CH3)2C=CH2 + HBF4 (CH3)3C+ +
The protonation of ammonia in the formation of ammonium chloride from ammonia and hydrogen chloride:
NH3(g) + HCl(g) → NH4Cl(s)
Protonation is a fundamental chemical reaction and is a step in many stoichiometric and catalytic processes. Some ions and molecules can undergo more than one protonation and are labeled polybasic, which is true of many biological macromolecules. Protonation and deprotonation (removal of a proton) occur in most acid–base reactions; they are the core of most acid–base reaction theories. A Brønsted–Lowry acid is defined as a chemical substance that protonates another substance. Upon protonating a substrate, the mass and the charge of the species each increase by one unit, making it an essential step in certain analytical procedures such as electrospray mass spectrometry. Protonating or deprotonating a molecule or ion can change many other chemical properties, not just the charge and mass, for example solubility, hydrophilicity, reduction potential or oxidation potential, and optical properties can change.
Rates
Protonations are often rapid, partly because of the high mobility of protons in many solvents. The rate of protonation is related to the acidity of the protonating species: protonation by weak acids is slower than protonation of the same base by strong acids. The rates of protonation and deprotonation can be especially slow when protonation induces significant structural changes.
Enantioselective protonations are under kinetic control, are of considerable interest in organic synthesis. They are also relevant to various biological processes.
Reversibility and catalysis
Protonation is usually reversible, and the structure and bonding of the conjugate base are normally unchanged on protonation. In some cases, however, protonation induces isomerization, for example cis-alkenes can be converted to trans-alkenes using a catalytic amount of protonating agent. Many enzymes, such as the serine hydrolases, operate by mechanisms that involve reversible protonation of substrates.
| Physical sciences | Concepts | Chemistry |
522274 | https://en.wikipedia.org/wiki/Deprotonation | Deprotonation | Deprotonation (or dehydronation) is the removal (transfer) of a proton (or hydron, or hydrogen cation), (H+) from a Brønsted–Lowry acid in an acid–base reaction. The species formed is the conjugate base of that acid. The complementary process, when a proton is added (transferred) to a Brønsted–Lowry base, is protonation (or hydronation). The species formed is the conjugate acid of that base.
A species that can either accept or donate a proton is referred to as amphiprotic. An example is the H2O (water) molecule, which can gain a proton to form the hydronium ion, H3O+, or lose a proton, leaving the hydroxide ion, OH−.
The relative ability of a molecule to give up a proton is measured by its pKa value. A low pKa value indicates that the compound is acidic and will easily give up its proton to a base. The pKa of a compound is determined by many aspects, but the most significant is the stability of the conjugate base. This is primarily determined by the ability (or inability) of the conjugated base to stabilize negative charge. One of the most important ways of assessing a conjugate base's ability to distribute negative charge is using resonance. Electron withdrawing groups (which can stabilize the molecule by increasing charge distribution) or electron donating groups (which destabilize by decreasing charge distribution) present on a molecule also determine its pKa. The solvent used can also assist in the stabilization of the negative charge on a conjugated base.
Bases used to deprotonate depend on the pKa of the compound. When the compound is not particularly acidic, and, as such, the molecule does not give up its proton easily, a base stronger than the commonly known hydroxides is required. Hydrides are one of the many types of powerful deprotonating agents. Common hydrides used are sodium hydride and potassium hydride. The hydride forms hydrogen gas with the liberated proton from the other molecule. The hydrogen is dangerous and could ignite with the oxygen in the air, so the chemical procedure should be done in an inert atmosphere (e.g., nitrogen).
Deprotonation can be an important step in a chemical reaction. Acid–base reactions typically occur faster than any other step which may determine the product of a reaction. The conjugate base is more electron-rich than the molecule which can alter the reactivity of the molecule. For example, deprotonation of an alcohol forms the negatively charged alkoxide, which is a much stronger nucleophile.
To determine whether or not a given base will be sufficient to deprotonate a specific acid, compare the conjugate base with the original base. A conjugate base is formed when the acid is deprotonated by the base. In the image above, hydroxide acts as a base to deprotonate the carboxylic acid. The conjugate base is the carboxylate salt. In this case, hydroxide is a strong enough base to deprotonate the carboxylic acid because the conjugate base is more stable than the base because the negative charge is delocalized over two electronegative atoms compared to one. Using pKa values, the carboxylic acid is approximately 4 and the conjugate acid, water, is 15.7. Because acids with higher pKa values are less likely to donate their protons, the equilibrium will favor their formation. Therefore, the side of the equation with water will be formed preferentially. If, for example, water, instead of hydroxide, was used to deprotonate the carboxylic acid, the equilibrium would not favor the formation of the carboxylate salt. This is because the conjugate acid, hydronium, has a pKa of -1.74, which is lower than the carboxylic acid. In this case, equilibrium would favor the carboxylic acid.
| Physical sciences | Concepts | Chemistry |
522347 | https://en.wikipedia.org/wiki/Pelagic%20zone | Pelagic zone | The pelagic zone consists of the water column of the open ocean and can be further divided into regions by depth. The word pelagic is derived . The pelagic zone can be thought of as an imaginary cylinder or water column between the surface of the sea and the bottom. Conditions in the water column change with depth: pressure increases; temperature and light decrease; salinity, oxygen, micronutrients (such as iron, magnesium and calcium) all change. In a manner analogous to stratification in the Earth's atmosphere, the water column can be divided vertically into up to five different layers (illustrated in the diagram), with the number of layers depending on the depth of the water.
Marine life is affected by bathymetry (underwater topography) such as the seafloor, shoreline, or a submarine seamount, as well as by proximity to the boundary between the ocean and the atmosphere at the ocean surface, which brings light for photosynthesis, predation from above, and wind stirring up waves and setting currents in motion. The pelagic zone refers to the open, free waters away from the shore, where marine life can swim freely in any direction unhindered by topographical constraints.
The oceanic zone is the deep open ocean beyond the continental shelf, which contrasts with the inshore waters near the coast, such as in estuaries or on the continental shelf. Waters in the oceanic zone plunge to the depths of the abyssopelagic and further to the hadopelagic. Coastal waters are generally the relatively shallow epipelagic. Altogether, the pelagic zone occupies 1,330 million km3 (320 million mi3) with a mean depth of and maximum depth of . Pelagic life decreases as depth increases.
The pelagic zone contrasts with the benthic and demersal zones at the bottom of the sea. The benthic zone is the ecological region at the very bottom, including the sediment surface and some subsurface layers. Marine organisms such as clams and crabs living in this zone are called benthos. Just above the benthic zone is the demersal zone. Demersal fish can be divided into benthic fish, which are denser than water and rest on the bottom, and benthopelagic fish, which swim just above the bottom. Demersal fish are also known as bottom feeders and groundfish.
Depth and layers
The pelagic zone is subdivided into five vertical regions. From the top down, these are:
Epipelagic (sunlight)
The illuminated zone at the surface of the sea with sufficient light for photosynthesis. Nearly all primary production in the ocean occurs here, and marine life is concentrated in this zone, including plankton, floating seaweed, jellyfish, tuna, many sharks and dolphins.
Mesopelagic (twilight)
The most abundant organisms thriving into the mesopelagic zone are heterotrophic bacteria. Animals living in this zone include swordfish, squid, wolffish and some species of cuttlefish. Many organisms living here are bioluminescent. Some mesopelagic creatures rise to the epipelagic zone at night to feed.
Bathypelagic (midnight)
The name stems . The ocean is pitch black at this depth apart from occasional bioluminescent organisms, such as anglerfish. No plants live here. Most animals survive on detritus known as "marine snow" falling from the zones above or, like the marine hatchetfish, by preying on other inhabitants of this zone. Other examples of this zone's inhabitants are giant squid, smaller squid and the grimpoteuthis or "dumbo octopus". The giant squid is hunted here by deep-diving sperm whales.
Abyssopelagic (abyssal zone)
The name is derived - a holdover from times when the deep ocean was believed to indeed be bottomless. Among the very few creatures living in the cold temperatures, high pressures and complete darkness here are several species of squid; echinoderms including the basket star, swimming cucumber, and the sea pig; and marine arthropods including the sea spider. Many species at these depths are transparent and eyeless.
Hadopelagic (hadal zone)
The name is derived from the realm of Hades, the Greek underworld. This is the deepest part of the ocean at more than or , depending on authority. Such depths are generally located in trenches.
Pelagic ecosystem
The pelagic ecosystem is based on phytoplankton. Phytoplankton manufacture their own food using a process of photosynthesis. Because they need sunlight, they inhabit the upper, sunlit epipelagic zone, which includes the coastal or neritic zone. Biodiversity diminishes markedly in the deeper zones below the epipelagic zone as dissolved oxygen diminishes, water pressure increases, temperatures become colder, food sources become scarce, and light diminishes and finally disappears.
Pelagic invertebrates
Some examples of pelagic invertebrates include krill, copepods, jellyfish, decapod larvae, hyperiid amphipods, rotifers and cladocerans.
Thorson's rule states that benthic marine invertebrates at low latitudes tend to produce large numbers of eggs developing to widely dispersing pelagic larvae, whereas at high latitudes such organisms tend to produce fewer and larger lecithotrophic (yolk-feeding) eggs and larger offspring.
Pelagic fish
Pelagic fish live in the water column of coastal, ocean, and lake waters, but not on or near the bottom of the sea or the lake. They can be contrasted with demersal fish, which do live on or near the bottom, and coral reef fish.
Pelagic fish are often migratory forage fish, which feed on plankton, and the larger predatory fish that follow and feed on the forage fish. Examples of migratory forage fish are herring, anchovies, capelin, and menhaden. Examples of larger pelagic fish which prey on the forage fish are billfish, tuna, and oceanic sharks.
Pelagic reptiles
Hydrophis platurus, the yellow-bellied sea snake, is the only one of the 65 species of marine snakes to spend its entire life in the pelagic zone. It bears live young at sea and is helpless on land. The species sometimes forms aggregations of thousands along slicks in surface waters. The yellow-bellied sea snake is the world's most widely distributed snake species.
Many species of sea turtles spend the first years of their lives in the pelagic zone, moving closer to shore as they reach maturity.
Pelagic birds
Pelagic birds, also called oceanic birds or seabirds, live on open seas and oceans rather than inland or around more restricted waters such as rivers and lakes. Pelagic birds feed on planktonic crustaceans, squid and forage fish. Examples are the Atlantic puffin, macaroni penguins, sooty terns, shearwaters, and Procellariiformes such as the albatross, Procellariidae and petrels.
| Physical sciences | Oceanography | Earth science |
522449 | https://en.wikipedia.org/wiki/Requirements%20analysis | Requirements analysis | In systems engineering and software engineering, requirements analysis focuses on the tasks that determine the needs or conditions to meet the new or altered product or project, taking account of the possibly conflicting requirements of the various stakeholders, analyzing, documenting, validating, and managing software or system requirements.
Requirements analysis is critical to the success or failure of a systems or software project. The requirements should be documented, actionable, measurable, testable, traceable, related to identified business needs or opportunities, and defined to a level of detail sufficient for system design.
Overview
Conceptually, requirements analysis includes three types of activities:
Eliciting requirements: (e.g. the project charter or definition), business process documentation, and stakeholder interviews. This is sometimes also called requirements gathering or requirements discovery.
Recording requirements: Requirements may be documented in various forms, usually including a summary list, and may include natural-language documents, use cases, user stories, process specifications, and a variety of models including data models.
Analyzing requirements: determining whether the stated requirements are clear, complete, unduplicated, concise, valid, consistent and unambiguous, and resolving any apparent conflicts. Analyzing can also include sizing requirements.
Requirements analysis can be a long and tiring process during which many delicate psychological skills are involved. New systems change the environment and relationships between people, so it is important to identify all the stakeholders, take into account all their needs, and ensure they understand the implications of the new systems. Analysts can employ several techniques to elicit the requirements from the customer. These may include the development of scenarios (represented as user stories in agile methods), the identification of use cases, the use of workplace observation or ethnography, holding interviews, or focus groups (more aptly named in this context as requirements workshops, or requirements review sessions) and creating requirements lists. Prototyping may be used to develop an example system that can be demonstrated to stakeholders. Where necessary, the analyst will employ a combination of these methods to establish the exact requirements of the stakeholders, so that a system that meets the business needs is produced. Requirements quality can be improved through these and other methods:
Visualization. Using tools that promote better understanding of the desired end-product such as visualization and simulation.
Consistent use of templates. Producing a consistent set of models and templates to document the requirements.
Documenting dependencies. Documenting dependencies and interrelationships among requirements, as well as any assumptions and congregations.
Requirements analysis topics
Stakeholder identification
See Stakeholder analysis for a discussion of people or organizations (legal entities such as companies, and standards bodies) that have a valid interest in the system. They may be affected by it either directly or indirectly.
A major new emphasis in the 1990s was a focus on the identification of stakeholders. It is increasingly recognized that stakeholders are not limited to the organization employing the analyst. Other stakeholders will include:
anyone who operates the system (normal and maintenance operators)
anyone who benefits from the system (functional, political, financial, and social beneficiaries)
anyone involved in purchasing or procuring the system. In a mass-market product organization, product management, marketing, and sometimes sales act as surrogate consumers (mass-market customers) to guide the development of the product.
organizations that regulate aspects of the system (financial, safety, and other regulators)
people or organizations opposed to the system (negative stakeholders; see also Misuse case)
organizations responsible for systems that interface with the system under design.
those organizations that integrate horizontally with the organization for whom the analyst is designing the system.
Joint Requirements Development (JRD) Sessions
Requirements often have cross-functional implications that are unknown to individual stakeholders and often missed or incompletely defined during stakeholder interviews. These cross-functional implications can be elicited by conducting JRD sessions in a controlled environment, facilitated by a trained facilitator (Business Analyst), wherein stakeholders participate in discussions to elicit requirements, analyze their details, and uncover cross-functional implications. A dedicated scribe should be present to document the discussion, freeing up the Business Analyst to lead the discussion in a direction that generates appropriate requirements that meet the session objective.
JRD Sessions are analogous to Joint Application Design Sessions. In the former, the sessions elicit requirements that guide design, whereas the latter elicit the specific design features to be implemented in satisfaction of elicited requirements.
Contract-style requirement lists
One traditional way of documenting requirements has been contract-style requirement lists. In a complex system such requirements lists can run hundreds of pages long.
An appropriate metaphor would be an extremely long shopping list. Such lists are very much out of favor in modern analysis; as they have proved spectacularly unsuccessful at achieving their aims; but they are still seen to this day.
Strengths
Provides a checklist of requirements.
Provide a contract between the project sponsor(s) and developers.
For a large system can provide a high level description from which lower-level requirements can be derived.
Weaknesses
Such lists can run to hundreds of pages. They are not intended to serve as a reader-friendly description of the desired application.
Such requirements lists abstract all the requirements and so there is little context. The Business Analyst may include context for requirements in accompanying design documentation.
This abstraction is not intended to describe how the requirements fit or work together.
The list may not reflect relationships and dependencies between requirements. While a list does make it easy to prioritize each item, removing one item out of context can render an entire use case or business requirement useless.
The list does not supplant the need to review requirements carefully with stakeholders to gain a better-shared understanding of the implications for the design of the desired system/application.
Simply creating a list does not guarantee its completeness. The Business Analyst must make a good faith effort to discover and collect a substantially comprehensive list and rely on stakeholders to point out missing requirements.
These lists can create a false sense of mutual understanding between the stakeholders and developers; Business Analysts are critical to the translation process.
It is almost impossible to uncover all the functional requirements before the process of development and testing begins. If these lists are treated as an immutable contract, then requirements that emerge in the Development process may generate a controversial change request.
Alternative to requirement lists
As an alternative to requirement lists, Agile Software Development uses User stories to suggest requirements in everyday language.
Measurable goals
Best practices take the composed list of requirements merely as clues and repeatedly ask "why?" until the actual business purposes are discovered. Stakeholders and developers can then devise tests to measure what level of each goal has been achieved thus far. Such goals change more slowly than the long list of specific but unmeasured requirements. Once a small set of critical, measured goals has been established, rapid prototyping and short iterative development phases may proceed to deliver actual stakeholder value long before the project is half over.
Prototypes
A prototype is a computer program that exhibits a part of the properties of another computer program, allowing users to visualize an application that has not yet been constructed. A popular form of prototype is a mockup, which helps future users and other stakeholders get an idea of what the system will look like. Prototypes make it easier to make design decisions because aspects of the application can be seen and shared before the application is built. Major improvements in communication between users and developers were often seen with the introduction of prototypes. Early views of applications led to fewer changes later and hence reduced overall costs considerably.
Prototypes can be flat diagrams (often referred to as wireframes) or working applications using synthesized functionality. Wireframes are made in a variety of graphic design documents, and often remove all color from the design (i.e. use a greyscale color palette) in instances where the final software is expected to have a graphic design applied to it. This helps to prevent confusion as to whether the prototype represents the final visual look and feel of the application.
Use cases
A use case is a structure for documenting the functional requirements for a system, usually involving software, whether that is new or being changed. Each use case provides a set of scenarios that convey how the system should interact with a human user or another system, to achieve a specific business goal. Use cases typically avoid technical jargon, preferring instead the language of the end-user or domain expert. Use cases are often co-authored by requirements engineers and stakeholders.
Use cases are deceptively simple tools for describing the behavior of software or systems. A use case contains a textual description of how users are intended to work with the software or system. Use cases should not describe the internal workings of the system, nor should they explain how that system will be implemented. Instead, they show the steps needed to perform a task without sequential assumptions.
Requirements specification
Requirements specification is the synthesis of discovery findings regarding current state business needs and the assessment of these needs to determine, and specify, what is required to meet the needs within the solution scope in focus. Discovery, analysis, and specification move the understanding from a current as-is state to a future to-be state. Requirements specification can cover the full breadth and depth of the future state to be realized, or it could target specific gaps to fill, such as priority software system bugs to fix and enhancements to make. Given that any large business process almost always employs software and data systems and technology, requirements specification is often associated with software system builds, purchases, cloud computing strategies, embedded software in products or devices, or other technologies. The broader definition of requirements specification includes or focuses on any solution strategy or component, such as training, documentation guides, personnel, marketing strategies, equipment, supplies, etc.
Types of requirements
Requirements are categorized in several ways. The following are common categorizations of requirements that relate to technical management:
Business requirements
Statements of business level goals, without reference to detailed functionality. These are usually high-level (software and/or hardware) capabilities that are needed to achieve a business outcome.
Customer requirements
Statements of fact and assumptions that define the expectations of the system in terms of mission objectives, environment, constraints, and measures of effectiveness and suitability (MOE/MOS). The customers are those that perform the eight primary functions of systems engineering, with special emphasis on the operator as the key customer. Operational requirements will define the basic need and, at a minimum, answer the questions posed in the following listing:
Operational distribution or deployment: Where will the system be used?
Mission profile or scenario: How will the system accomplish its mission objective?
Performance and related parameters: What are the critical system parameters to accomplish the mission?
Utilization environments: How are the various system components to be used?
Effectiveness requirements: How effective or efficient must the system be in performing its mission?
Operational life cycle: How long will the system be in use by the user?
Environment: What environments will the system be expected to operate in an effective manner?
Architectural requirements
Architectural requirements explain what has to be done by identifying the necessary systems architecture of a system.
Behavioral requirements
Behavioral requirements explain what has to be done by identifying the necessary behavior of a system.
Functional requirements
Functional requirements explain what has to be done by identifying the necessary task, action or activity that must be accomplished. Functional requirements analysis will be used as the toplevel functions for functional analysis.
Non-functional requirements
Non-functional requirements are requirements that specify criteria that can be used to judge the operation of a system, rather than specific behaviors.
Performance requirements
The extent to which a mission or function must be executed; is generally measured in terms of quantity, quality, coverage, timeliness, or readiness. During requirements analysis, performance (how well does it have to be done) requirements will be interactively developed across all identified functions based on system life cycle factors; and characterized in terms of the degree of certainty in their estimate, the degree of criticality to the system success, and their relationship to other requirements.
Design requirements
The "build to", "code to", and "buy to" requirements for products and "how to execute" requirements for processes are expressed in technical data packages and technical manuals.
Derived requirements
Requirements that are implied or transformed from higher-level requirements. For example, a requirement for long-range or high speed may result in a design requirement for low weight.
Allocated requirements
A requirement is established by dividing or otherwise allocating a high-level requirement into multiple lower-level requirements. Example: A 100-pound item that consists of two subsystems might result in weight requirements of 70 pounds and 30 pounds for the two lower-level items.
Well-known requirements categorization models include FURPS and FURPS+, developed at Hewlett-Packard.
Requirements analysis issues
Stakeholder issues
Steve McConnell, in his book Rapid Development, details a number of ways users can inhibit requirements gathering:
Users do not understand what they want or users do not have a clear idea of their requirements
Users will not commit to a set of written requirements
Users insist on new requirements after the cost and schedule have been fixed
Communication with users is slow
Users often do not participate in reviews or are incapable of doing so
Users are technically unsophisticated
Users do not understand the development process
Users do not know about present technology
This may lead to the situation where user requirements keep changing even when system or product development has been started.
Engineer/developer issues
Possible problems caused by engineers and developers during requirements analysis are:
A natural inclination towards writing code can lead to implementation beginning before the requirements analysis is complete, potentially resulting in code changes to meet actual requirements once they are known.
Technical personnel and end-users may have different vocabularies. Consequently, they may wrongly believe they are in perfect agreement until the finished product is supplied.
Engineers and developers may try to make the requirements fit an existing system or model, rather than develop a system specific to the needs of the client.
Attempted solutions
One attempted solution to communications problems has been to employ specialists in business or system analysis.
Techniques introduced in the 1990s like prototyping, Unified Modeling Language (UML), use cases, and agile software development are also intended as solutions to problems encountered with previous methods.
Also, a new class of application simulation or application definition tools has entered the market. These tools are designed to bridge the communication gap between business users and the IT organization — and also to allow applications to be 'test marketed' before any code is produced. The best of these tools offer:
electronic whiteboards to sketch application flows and test alternatives
ability to capture business logic and data needs
ability to generate high-fidelity prototypes that closely imitate the final application
interactivity
capability to add contextual requirements and other comments
ability for remote and distributed users to run and interact with the simulation
| Technology | Basics | null |
522476 | https://en.wikipedia.org/wiki/Stereocenter | Stereocenter | In stereochemistry, a stereocenter of a molecule is an atom (center), axis or plane that is the focus of stereoisomerism; that is, when having at least three different groups bound to the stereocenter, interchanging any two different groups creates a new stereoisomer. Stereocenters are also referred to as stereogenic centers.
A stereocenter is geometrically defined as a point (location) in a molecule; a stereocenter is usually but not always a specific atom, often carbon. Stereocenters can exist on chiral or achiral molecules; stereocenters can contain single bonds or double bonds. The number of hypothetical stereoisomers can be predicted by using 2n, with n being the number of tetrahedral stereocenters; however, exceptions such as meso compounds can reduce the prediction to below the expected 2n.
Chirality centers are a type of stereocenter with four different substituent groups; chirality centers are a specific subset of stereocenters because they can only have sp3 hybridization, meaning that they can only have single bonds.
Location
Stereocenters can exist on chiral or achiral molecules. They are defined as a location (point) within a molecule, rather than a particular atom, in which the interchanging of two groups creates a stereoisomer. A stereocenter can have either four different attachment groups, or three different attachment groups where one group is connected by a double bond. Since stereocenters can exist on achiral molecules, stereocenters can have either sp3 or sp2 hybridization.
Possible number of stereoisomers
Stereoisomers are compounds that are identical in composition and connectivity but have a different spatial arrangement of atoms around the central atom. A molecule having multiple stereocenters will produce many possible stereoisomers. In compounds whose stereoisomerism is due to tetrahedral (sp3) stereogenic centers, the total number of hypothetically possible stereoisomers will not exceed 2n, where n is the number of tetrahedral stereocenters. However, this is an upper bound because molecules with symmetry frequently have fewer stereoisomers.
The stereoisomers produced by the presence of multiple stereocenters can be defined as enantiomers (non-superposable mirror images) and diastereomers (non-superposable, non-identical, non-mirror image molecules). Enantiomers and diastereomers are produced due to differing stereochemical configurations of molecules containing the same composition and connectivity (bonding); the molecules must have multiple (two or more) stereocenters to be classified as enantiomers or diastereomers. Enantiomers and diastereomers will produce individual stereoisomers that contribute to the total number of possible stereoisomers.
However, the stereoisomers produced may also give a meso compound, which is an achiral compound that is superposable on its mirror image; the presence of a meso compound will reduce the number of possible stereoisomers. Since a meso compound is superposable on its mirror image, the two "stereoisomers" are actually identical. Resultantly, a meso compound will reduce the number of stereoisomers to below the hypothetical 2n amount due to symmetry.
Additionally, certain configurations may not exist due to steric reasons. Cyclic compounds with chiral centers may not exhibit chirality due to the presence of a two-fold rotation axis. Planar chirality may also provide for chirality without having an actual chiral center present.
Configuration
Configuration is defined as the arrangement of atoms around a stereocenter. The Cahn-Ingold-Prelog (CIP) system uses R and S designations to define the configuration of atoms about any stereocenter. A designation of R denotes a clockwise direction of substituent priority around the stereocenter, while a designation of S denotes a counter-clockwise direction of substituent priority.
Chirality centers
A chirality center (chiral center) is a type of stereocenter. A chirality center is defined as an atom holding a set of four different ligands (atoms or groups of atoms) in a spatial arrangement which is non-superposable on its mirror image. Chirality centers must be sp3 hybridized, meaning that a chirality center can only have single bonds. In organic chemistry, a chirality center usually refers to a carbon, phosphorus, or sulfur atom, though it is also possible for other atoms to be chirality centers, especially in areas of organometallic and inorganic chemistry.
The concept of a chirality center generalizes the concept of an asymmetric carbon atom (a carbon atom bonded to four different entities) to a broader definition of any atom with four different attachment groups in which an interchanging of any two attachment groups gives rise to an enantiomer.
Stereogenic on carbon
A carbon atom that is attached to four different substituent groups is called an asymmetric carbon atom or chiral carbon. Chiral carbons are the most common type of chirality center.
Stereogenic on other atoms
Chirality is not limited to carbon atoms, though carbon atoms are often centers of chirality due to their ubiquity in organic chemistry. Nitrogen and phosphorus atoms can also form bonds in a tetrahedral configuration. A nitrogen in an amine may be a stereocenter if all three groups attached are different because the electron pair of the amine functions as a fourth group. However, nitrogen inversion, a form of pyramidal inversion, causes racemization which means that both epimers at that nitrogen are present under normal circumstances. Racemization by nitrogen inversion may be restricted (such as quaternary ammonium or phosphonium cations), or slow, which allows the existence of chirality.
Metal atoms with tetrahedral or octahedral geometries may also be chiral due to having different ligands. For the octahedral case, several chiralities are possible. Having three ligands of two types, the ligands may be lined up along the meridian, giving the mer-isomer, or forming a face—the fac isomer. Having three bidentate ligands of only one type gives a propeller-type structure, with two different enantiomers denoted Λ and Δ.
Chirality and stereocenters
As mentioned earlier, the requirement for an atom to be a chirality center is that the atom must be sp3 hybridized with four different attachments. Because of this, all chirality centers are stereocenters. However, only under some conditions is the reverse true. Recall that a point can be considered a sterocenter with a minimum of three attachment points; stereocenters can be either sp3 or sp2 hybridized, as long as the interchanging any two different groups creates a new stereoisomer. This means that although all chirality centers are stereocenters, not every stereocenter is a chirality center.
Stereocenters are important identifiers for chiral or achiral molecules. As a general rule, if a molecule has no stereocenters, it is considered achiral. If it has at least one stereocenter, the molecule has the potential for chirality. However, there are some exceptions like meso compounds that make molecules with multiple stereocenters considered achiral.
| Physical sciences | Stereochemistry | Chemistry |
522519 | https://en.wikipedia.org/wiki/Diastereomer | Diastereomer | In stereochemistry, diastereomers (sometimes called diastereoisomers) are a type of stereoisomer. Diastereomers are defined as non-mirror image, non-identical stereoisomers. Hence, they occur when two or more stereoisomers of a compound have different configurations at one or more (but not all) of the equivalent (related) stereocenters and are not mirror images of each other.
When two diastereoisomers differ from each other at only one stereocenter, they are epimers. Each stereocenter gives rise to two different configurations and thus typically increases the number of stereoisomers by a factor of two.
Diastereomers differ from enantiomers in that the latter are pairs of stereoisomers that differ in all stereocenters and are therefore mirror images of one another.
Enantiomers of a compound with more than one stereocenter are also diastereomers of the other stereoisomers of that compound that are not their mirror image (that is, excluding the opposing enantiomer).
Diastereomers have different physical properties (unlike most aspects of enantiomers) and often different chemical reactivity.
Diastereomers differ not only in physical properties but also in chemical reactivity — how a compound reacts with others. Glucose and galactose, for instance, are diastereomers. Even though they share the same molar weight, glucose is more stable than galactose. This difference in stability causes galactose to be absorbed slightly faster than glucose in human body.
Diastereoselectivity is the preference for the formation of one or more than one diastereomer over the other in an organic reaction. In general, stereoselectivity is attributed to torsional and steric interactions in the stereocenter resulting from electrophiles approaching the stereocenter in reaction.
Syn / anti
When the single bond between the two centres is free to rotate, cis/trans descriptors become invalid. Two widely accepted prefixes used to distinguish diastereomers on sp³-hybridised bonds in an open-chain molecule are syn and anti. Masamune proposed the descriptors which work even if the groups are not attached to adjacent carbon atoms. It also works regardless of CIP priorities. Syn describes groups on the same face while anti describes groups on opposite faces. The concept applies only to the Zigzag projection. The descriptors only describe relative stereochemistry rather than absolute stereochemistry.
All isomers are same.
Erythro / threo
Two older prefixes still commonly used to distinguish diastereomers are threo and erythro. In the case of saccharides, when drawn in the Fischer projection the erythro isomer has two identical substituents on the same side and the threo isomer has them on opposite sides. When drawn as a zig-zag chain, the erythro isomer has two identical substituents on different sides of the plane (anti). The names are derived from the diastereomeric four-carbon aldoses erythrose and threose. These prefixes are not recommended for general use because it is often difficult to discern how to apply their definitions to particular comounds. However, the prefixes can usefully describe the relative configuration of a compound that has the following properties: it has at least four C atoms, exactly two of those C atoms are stereocenters, the stereocenters are adjacent, and the two substituents on each stereocenter can clearly be labeled as "larger" (usually a heteroatom such as N, O, or S) and "smaller" (usually H).
Threitol and erythritol are both four-carbon sugar alcohols. Erythritol is achiral (has at least one conformation with a plane or center of symmetry), whereas threitol is chiral. A useful English-language mnemonic device is that "threitol" and "chiral" both begin with consonants, whereas "erythritol" and "achiral" both begin with vowels.
Another threo compound is threonine, one of the amino acids coded by DNA. Its erythro diastereomer, allothreonine, is not coded by DNA and is very rare in nature.
In alkene addition reactions, syn addition to a trans alkene, or anti addition to a cis alkene, gives a threo product, whereas syn addition to a cis alkene, or anti addition to a trans alkene, gives an erythro product.
Multiple stereocenters
If a molecule contains two asymmetric centers, there are up to four possible configurations, and they cannot all be non-superposable mirror images of each other. The possibilities for different isomers continue to multiply as more stereocenters are added to a molecule. In general, the number of stereoisomers of a molecule can be determined by calculating 2n, where n = the number of chiral centers in the molecule. This holds true except in cases where the molecule has meso forms. These meso compounds are molecules that contain stereocenters, but possess an internal plane of symmetry allowing it to be superposed on its mirror image. These equivalent configurations cannot be considered diastereomers.
For n = 3, there are eight stereoisomers. Among them, there are four pairs of enantiomers: R,R,R and S,S,S; R,R,S and S,S,R; R,S,S and S,R,R; and R,S,R and S,R,S. There are many more pairs of diastereomers, because each of these configurations is a diastereomer with respect to every other configuration excluding its own enantiomer (for example, R,R,R is a diastereomer of R,R,S; R,S,R; and R,S,S). For n = 4, there are sixteen stereoisomers, or eight pairs of enantiomers. The four enantiomeric pairs of aldopentoses and the eight enantiomeric pairs of aldohexoses (subsets of the five- and six-carbon sugars) are examples of sets of compounds that differ in this way.
Diastereomerism at a double bond
Double bond isomers are always considered diastereomers, not enantiomers. Diastereomerism can also occur at a double bond, where the cis vs trans relative positions of substituents give two non-superposable isomers. Many conformational isomers are diastereomers as well.
In the case of diastereomerism occurring at a double bond, E-Z, or entgegen and zusammen (German), is used in notating nomenclature of alkenes.
Applications
As stated previously, two diastereomers will not have identical chemical properties. This knowledge is harnessed in chiral synthesis to separate a mixture of enantiomers. This is the principle behind chiral resolution. After preparing the diastereomers, they are separated by chromatography or recrystallization. Note also the example of the stereochemistry of ketonization of enols and enolates.
| Physical sciences | Stereochemistry | Chemistry |
522690 | https://en.wikipedia.org/wiki/Oxalic%20acid | Oxalic acid | Oxalic acid is an organic acid with the systematic name ethanedioic acid and chemical formula , also written as or or . It is the simplest dicarboxylic acid. It is a white crystalline solid that forms a colorless solution in water. Its name comes from the fact that early investigators isolated oxalic acid from flowering plants of the genus Oxalis, commonly known as wood-sorrels. It occurs naturally in many foods. Excessive ingestion of oxalic acid or prolonged skin contact can be dangerous.
Oxalic acid has much greater acid strength than acetic acid. It is a reducing agent and its conjugate bases hydrogen oxalate () and oxalate () are chelating agents for metal cations. It is used as a cleaning agent, especially for the removal of rust, because it forms a water-soluble ferric iron complex, the ferrioxalate ion. Oxalic acid typically occurs as the dihydrate with the formula .
History
The preparation of salts of oxalic acid from plants had been known, at least since 1745, when the Dutch botanist and physician Herman Boerhaave isolated a salt from wood sorrel, akin to kraft process. By 1773, François Pierre Savary of Fribourg, Switzerland had isolated oxalic acid from its salt in sorrel.
In 1776, Swedish chemists Carl Wilhelm Scheele and Torbern Olof Bergman produced oxalic acid by reacting sugar with concentrated nitric acid; Scheele called the acid that resulted socker-syra or såcker-syra (sugar acid). By 1784, Scheele had shown that "sugar acid" and oxalic acid from natural sources were identical. The modern name was introduced along with many other acid names by de Morveau, Lavoisier and coauthors in 1787.
In 1824, the German chemist Friedrich Wöhler obtained oxalic acid by reacting cyanogen with ammonia in aqueous solution. This experiment may represent the first synthesis of a natural product.
Production
Industrial
Oxalic acid is mainly manufactured by the oxidation of carbohydrates or glucose using nitric acid or air in the presence of vanadium pentoxide. Another process uses oxygen to regenerate the nitric acid, using a variety of precursors including glycolic acid and ethylene glycol. As of 2011, this process was only used by Mitsubishi in Japan. A newer method entails oxidative carbonylation of alcohols to give the diesters of oxalic acid:
These diesters are subsequently hydrolyzed to oxalic acid. Approximately 120,000 tonnes are produced annually.
Historically oxalic acid was obtained exclusively by using caustics, such as sodium or potassium hydroxide, on sawdust, followed by acidification of the oxalate by mineral acids, such as sulfuric acid. Oxalic acid can also be formed by the heating of sodium formate in the presence of an alkaline catalyst.
Laboratory
Although it can be readily purchased, oxalic acid can be prepared in the laboratory by oxidizing sucrose using nitric acid in the presence of a small amount of vanadium pentoxide as a catalyst.
The hydrated solid can be dehydrated with heat or by azeotropic distillation.
Structure
Anhydrous
Anhydrous oxalic acid exists as two polymorphs; in one the hydrogen-bonding results in a chain-like structure, whereas the hydrogen bonding pattern in the other form defines a sheet-like structure. Because the anhydrous material is both acidic and hydrophilic (water seeking), it is used in esterifications.
Dihydrate
The dihydrate ·2 has space group C52h–P21/n, with lattice parameters , , , , . The main inter-atomic distances are: C−C 153 pm, C−O1 129 pm, C−O2 119 pm.
Reactions
Acid–base properties
Oxalic acid's pKa values vary in the literature from 1.25 to 1.46 and from 3.81 to 4.40. The 100th ed of the CRC, released in 2019, has values of 1.25 and 3.81.
Oxalic acid is relatively strong compared to other carboxylic acids:
Oxalic acid undergoes many of the reactions characteristic for other carboxylic acids. It forms esters such as dimethyl oxalate (m.p. ). It forms an acid chloride called oxalyl chloride.
Metal-binding properties
Transition metal oxalate complexes are numerous, e.g. the drug oxaliplatin. Oxalic acid has been shown to reduce manganese dioxide in manganese ores to allow the leaching of the metal by sulfuric acid.
Oxalic acid is an important reagent in lanthanide chemistry. Hydrated lanthanide oxalates form readily in very strongly acidic solutions as a densely crystalline, easily filtered form, largely free of contamination by nonlanthanide elements:
Thermal decomposition of these oxalates gives the oxides, which is the most commonly marketed form of these elements.
Other
Oxalic acid and oxalates can be oxidized by permanganate in an autocatalytic reaction.
Oxalic acid vapor decomposes at 125–175 °C into carbon dioxide and formic acid HCOOH. Photolysis with 237–313 nm UV light also produces carbon monoxide CO and water.
Evaporation of a solution of urea and oxalic acid in 2:1 molar ratio yields a solid crystalline compound , consisting of stacked two-dimensional networks of the neutral molecules held together by hydrogen bonds with the oxygen atoms.
Occurrence
Biosynthesis
At least two pathways exist for the enzyme-mediated formation of oxalate. In one pathway, oxaloacetate, a component of the Krebs citric acid cycle, is hydrolyzed to oxalate and acetic acid by the enzyme oxaloacetase:
It also arises from the dehydrogenation of glycolic acid, which is produced by the metabolism of ethylene glycol.
Occurrence in foods and plants
Early investigators isolated oxalic acid from wood-sorrel (Oxalis). Members of the spinach family and the brassicas (cabbage, broccoli, brussels sprouts) are high in oxalates, as are sorrel and umbellifers like parsley. The leaves and stems of all species of the genus Chenopodium and related genera of the family Amaranthaceae, which includes quinoa, contain high levels of oxalic acid. Rhubarb leaves contain about 0.5% oxalic acid, and jack-in-the-pulpit (Arisaema triphyllum) contains calcium oxalate crystals. Similarly, the Virginia creeper, a common decorative vine, produces oxalic acid in its berries as well as oxalate crystals in the sap, in the form of raphides. Bacteria produce oxalates from oxidation of carbohydrates.
Plants of the genus Fenestraria produce optical fibers made from crystalline oxalic acid to transmit light to subterranean photosynthetic sites.
Carambola, also known as starfruit, also contains oxalic acid along with caramboxin. Citrus juice contains small amounts of oxalic acid.
The formation of naturally occurring calcium oxalate patinas on certain limestone and marble statues and monuments has been proposed to be caused by the chemical reaction of the carbonate stone with oxalic acid secreted by lichen or other microorganisms.
Production by fungi
Many soil fungus species secrete oxalic acid, which results in greater solubility of metal cations and increased availability of certain soil nutrients, and can lead to the formation of calcium oxalate crystals. Some fungi such as Aspergillus niger have been extensively studied for the industrial production of oxalic acid; however, those processes are not yet economically competitive with production from oil and gas.
Cryphonectria parasitica may excrete oxalic acid containing solutions at the advancing edge of its chestnut cambium infection. The lower pH (<2.5) of more concentrated oxalic acid excretions may degrade cambium cell walls and have a toxic effect on chestnut cambium cells. Cambium cells that burst provide nutrients for a blight infection advance.
Biochemistry
The conjugate base of oxalic acid is the hydrogenoxalate anion, and its conjugate base (oxalate) is a competitive inhibitor of the lactate dehydrogenase (LDH) enzyme. LDH catalyses the conversion of pyruvate to lactic acid (end product of the fermentation (anaerobic) process) oxidising the coenzyme NADH to NAD+ and H+ concurrently. Restoring NAD+ levels is essential to the continuation of anaerobic energy metabolism through glycolysis. As cancer cells preferentially use anaerobic metabolism (see Warburg effect) inhibition of LDH has been shown to inhibit tumor formation and growth, thus is an interesting potential course of cancer treatment.
Oxalic acid plays a key role in the interaction between pathogenic fungi and plants. Small amounts of oxalic acid enhances plant resistance to fungi, but higher amounts cause widespread programmed cell death of the plant and help with fungi infection. Plants normally produce it in small amounts, but some pathogenic fungi such as Sclerotinia sclerotiorum cause a toxic accumulation.
Oxalate, besides being biosynthesised, may also be biodegraded. Oxalobacter formigenes is an important gut bacterium that helps animals (including humans) degrade oxalate.
Applications
Oxalic acid's main applications include cleaning or bleaching, especially for the removal of rust (iron complexing agent). Its utility in rust removal agents is due to its forming a stable, water-soluble salt with ferric iron, ferrioxalate ion. Oxalic acid is an ingredient in some tooth whitening products. About 25% of produced oxalic acid is used as a mordant in dyeing processes. It is also used in bleaches, especially for pulpwood, cork, straw, cane, feathers, and for rust removal and other cleaning, in baking powder, and as a third reagent in silica analysis instruments.
Niche uses
Oxalic acid is used by some beekeepers as a miticide against the parasitic varroa mite.
Dilute solutions (0.05–0.15 M) of oxalic acid can be used to remove iron from clays such as kaolinite to produce light-colored ceramics.
Oxalic acid can be used to clean minerals like many other acids. Two such examples are quartz crystals and pyrite.
Oxalic acid is sometimes used in the aluminum anodizing process, with or without sulfuric acid. Compared to sulfuric-acid anodizing, the coatings obtained are thinner and exhibit lower surface roughness.
Oxalic acid is also widely used as a wood bleach, most often in its crystalline form to be mixed with water to its proper dilution for use.
Semiconductor industry
Oxalic acid is also used in electronic and semiconductor industries. In 2006 it was reported being used in electrochemical–mechanical planarization of copper layers in the semiconductor devices fabrication process.
Proposed uses
Reduction of carbon dioxide to oxalic acid by various methods, such as electrocatalysis using a copper complex, is under study as a proposed chemical intermediate for carbon capture and utilization.
Content in food items
Toxicity
Oxalic acid has an oral LDLo (lowest published lethal dose) of 600 mg/kg. It has been reported that the lethal oral dose is 15 to 30 grams. The toxicity of oxalic acid is due to kidney failure caused by precipitation of solid calcium oxalate.
Oxalate is known to cause mitochondrial dysfunction.
Ingestion of ethylene glycol results in oxalic acid as a metabolite which can also cause acute kidney failure.
Kidney stones
Most kidney stones, 76%, are composed of calcium oxalate.
| Physical sciences | Specific acids | Chemistry |
522835 | https://en.wikipedia.org/wiki/Naegleriasis | Naegleriasis | Naegleriasis, also known as primary amoebic meningoencephalitis (PAM), is an almost invariably fatal infection of the brain by the free-living unicellular eukaryote Naegleria fowleri. Symptoms are meningitis-like and include headache, fever, nausea, vomiting, a stiff neck, confusion, hallucinations and seizures. Symptoms progress rapidly over around five days, and death usually results within one to two weeks of symptom onset.
N. fowleri is typically found in warm bodies of fresh water, such as ponds, lakes, rivers and hot springs. It is found in an amoeboid, temporary flagellate stage or microbial cyst in soil, poorly maintained municipal water supplies, water heaters, near warm-water discharges of industrial plants and in poorly chlorinated or unchlorinated swimming pools. There is no evidence of it living in salt water. As the disease is rare, it is often not considered during diagnosis.
Although infection occurs very rarely, it almost inevitably results in death. Of the 128 known USA naegleriasis cases in the half-century to 2016, only two survived.
Signs and symptoms
Onset of symptoms begins one to twelve days following exposure (with a median of five). Initial symptoms include changes in taste and smell, headache, fever, nausea, vomiting, back pain, and a stiff neck. Secondary symptoms are also meningitis-like including confusion, hallucinations, lack of attention, ataxia, cramp and seizures. After the start of symptoms, the disease progresses rapidly, with death usually occurring anywhere from one to eighteen days later (with a median of five), although it can take longer. In 2013, a man in Taiwan died 25 days after being infected by Naegleria fowleri.
It affects healthy children or young adults who have recently been exposed to bodies of fresh water. Scientists speculate that lower age groups are at a higher risk of contracting the disease because adolescents have a more underdeveloped and porous cribriform plate, through which the amoeba travels to reach the brain.
Cause
N. fowleri invades the central nervous system via the nose, specifically through the olfactory mucosa of the nasal tissues. This usually occurs as the result of the introduction of water that has been contaminated with N. fowleri into the nose during activities such as swimming, bathing or nasal irrigation.
The amoeba follows the olfactory nerve fibers through the cribriform plate of the ethmoid bone into the skull. There, it migrates to the olfactory bulbs and subsequently other regions of the brain, where it feeds on the nerve tissue. The organism then begins to consume cells of the brain, piecemeal through trogocytosis, by means of an amoebostome, a unique actin-rich sucking apparatus extended from its cell surface. It then becomes pathogenic, causing primary amoebic meningoencephalitis (PAM or PAME).
Primary amoebic meningoencephalitis presents symptoms similar to those of relatively common bacterial and viral meningitis. Upon abrupt disease onset, a plethora of symptoms arise. Endogenous cytokines, released in response to the pathogens, affect the thermoregulatory neurons of the hypothalamus causing a rise in body temperature. Additionally, the cytokines may act on the vascular organ of the lamina terminalis, leading to upregulation of Prostaglandin E2 contributing to hyperthermia. Further, the release of cytokines, exotoxins released by the pathogens and an increase in intracranial pressure stimulate the nociceptors in the meninges resulting in pain sensations.
The release of cytotoxic molecules in the central nervous system leads to extensive tissue damage and necrosis, such as damage to the olfactory nerve through lysis of nerve cells and demyelination. Specifically, the olfactory nerve and bulbs become necrotic and hemorrhagic. Spinal flexion leads to nuchal rigidity, or stiff neck, due to the stretching of the inflamed meninges. The increase in intracranial pressure stimulates the area postrema to create nausea sensations which may lead to brain herniation and damage to the reticular formation. Ultimately, the increase in cerebrospinal fluid from inflammation of the meninges increases intracranial pressure to an extent which leads to the destruction of the central nervous system. Although the exact pathophysiology behind the seizures caused by PAM is unknown, it is speculated that the seizures arise from altered meningeal permeability caused by increased intracranial pressure.
Pathogenesis
Naegleria fowleri propagates in warm, stagnant bodies of fresh water (typically during the summer months), and enters the central nervous system after insufflation of infected water by attaching itself to the olfactory nerve. It then migrates through the cribriform plate and into the olfactory bulbs of the forebrain, where it rapidly multiplies by feeding on nerve tissue.
Diagnosis
N. fowleri can be grown in several kinds of liquid axenic media or on non-nutrient agar plates coated with bacteria. Escherichia coli can be used to overlay the non-nutrient agar plate and a drop of cerebrospinal fluid sediment is added to it. Plates are then incubated at 37 °C and checked daily for clearing of the agar in thin tracks, which indicate the trophozoites have fed on the bacteria.
Detection in water is performed by centrifuging a water sample with E. coli added, then applying the pellet to a non-nutrient agar plate. After several days, the plate is microscopically inspected and Naegleria cysts are identified by their morphology. Final confirmation of the species' identity can be performed by various molecular or biochemical methods.
Confirmation of Naegleria presence can be done by a so-called flagellation test, where the organism is exposed to a hypotonic environment (distilled water). Naegleria, in contrast to other amoebae, differentiates within two hours into the flagellate state. Pathogenicity can be further confirmed by exposure to high temperature (42 °C): Naegleria fowleri is able to grow at this temperature, but the nonpathogenic Naegleria gruberi is not.
Prevention
Michael Beach, a recreational waterborne illness specialist for the Centers for Disease Control and Prevention, stated in remarks to the Associated Press that wearing of nose clips to prevent insufflation of contaminated water would be effective protection against contracting PAM, noting that "You'd have to have water going way up in your nose to begin with".
Advice stated in the press release from Taiwan's Centers for Disease Control recommended people prevent fresh water from entering the nostrils and avoid putting their heads down into fresh water or stirring mud in the water with feet. When starting to suffer from fever, headache, nausea, or vomiting subsequent to any kind of exposure to fresh water, even in the belief that no fresh water has traveled through the nostrils, people with such conditions should be carried to hospital quickly and make sure doctors are well-informed about the history of exposure to fresh water.
Treatment
On the basis of laboratory evidence and case reports, heroic doses of amphotericin B have been the traditional mainstay of PAM treatment since the first reported survivor in the United States in 1982.
Treatment has often also used combination therapy with multiple other antimicrobials in addition to amphotericin, such as fluconazole, miconazole, rifampicin and azithromycin. They have shown limited success only when administered early in the course of an infection.
While the use of rifampicin has been common, including in all four North American cases of survival, its continued use has been questioned. It only has variable activity in vitro and it has strong effects on the therapeutic levels of other antimicrobials used by inducing cytochrome p450 pathways. Fluconazole is commonly used as it has been shown to have synergistic effects against naegleria when used with amphotericin in vitro.
In 2013–2016, three successfully treated cases in the United States utilized the medication miltefosine. In one of the cases, a 12-year-old female, was given miltefosine and targeted temperature management to manage cerebral edema which is secondary to the infection. She survived with no neurological damage. The targeted temperature management coupled with early diagnosis and the medication has been attributed with her survival. On the other hand, another survivor, an 8-year-old male, was diagnosed several days after symptoms appeared and was not treated with targeted temperature management although he was administered miltefosine. He suffered apparent permanent neurological damage. In 2016, a 16-year-old male also survived PAM. He was treated with the same protocols as of the 12-year-old female in 2013. He recovered with a near-complete neurological recovery; however, the patient has mentioned difficulties with learning post-recovery.
the U.S. CDC offered miltefosine to doctors for the treatment of diseases caused by free-living amoebas including Naegleria, despite a lack of any data on how well the drug reaches the central nervous system.
In 2018, a 10-year-old girl in the Spanish city of Toledo became the first person to contract the disease in Spain, and was successfully treated using intravenous and intrathecal amphotericin B.
A 2023 study on mice has showed that treatment that included a derivative of the drug acoziborole known as AN3057 significantly prolonged survival and showed a 28% recovery rate without relapse.
Prognosis
Since its first description in the 1960s, only seven people worldwide have been reported to have survived PAM out of 450 cases diagnosed, implying a fatality rate of about 98.5%. The survivors include four in the United States, one in Mexico and one in Spain. One of the US survivors had brain damage that is likely permanent, but there are two documented surviving cases in the United States who made a full recovery with no neurological damage; they were both treated with the same protocols.
There is also a fourth survivor in the United States. However, he had a different strain.
Epidemiology
The disease is rare and highly lethal: there had only been 381 cases Drug treatment research at Aga Khan University in Pakistan has shown that in vitro drug susceptibility tests with some FDA approved drugs used for non-infectious diseases (digoxin and procyclidine were shown to be most effective of the drugs studied) have proved to kill Naegleria fowleri with an amoebicidal rate greater than 95%. The same source has also proposed a device for drug delivery via the transcranial route to the brain.
In the US, the most common states with cases reported of PAM from N. fowleri are the southern states, with Texas and Florida having the highest prevalence. The most commonly affected age group is 5–14-year olds (those who play in water). The number of cases of infection could increase due to climate change, which was posited as the reason for three cases in Minnesota in 2010, 2012, and 2015.
the numbers of reported cases were expected to increase simply because of better-informed diagnoses being made both in ongoing cases and in autopsy findings.
History
In 1899, Franz Schardinger first discovered and documented an amoeba he called Amoeba gruberi that could transform into a flagellate. The genus Naegleria was established by Alexis Alexeieff in 1912, who grouped the flagellate amoeba. He coined the term Naegleria after Kurt Nägler, who researched amoebae. It was not until 1965 that doctors Malcolm Fowler and Rodney F. Carter in Adelaide, Australia, reported the first four-human cases of amoebic meningoencephalitis. These cases involved four Australian children, one in 1961 and the rest in 1965, all of whom had succumbed to the illness. Their work on amebo-flagellates has provided an example of how a protozoan can effectively live both freely in the environment, and in a human host.
In 1966, Butt termed the infection resulting from N. fowleri primary amoebic meningoencephalitis (PAM) to distinguish this central nervous system (CNS) invasion from other secondary invasions made by other amoebae such as Entamoeba histolytica. A retrospective study determined the first documented case of PAM possibly occurred in Britain in 1909. In 1966, four cases were reported in the US. By 1968 the causative organism, previously thought to be a species of Hartmannella, was identified as a novel species of Naegleria. This same year, occurrence of sixteen cases over a period of three years (1962–1965) was reported in Ústí nad Labem, Czechoslovakia. In 1970, Carter named the species of amoeba N. fowleri, after Malcolm Fowler.
Society and culture
Naegleria fowleri is also known as the "brain-eating amoeba". This common name has also been applied to Balamuthia mandrillaris, causing some confusion between the two; Balamuthia mandrillaris is unrelated to Naegleria fowleri, and causes a different disease called granulomatous amoebic encephalitis. Unlike naegleriasis, which is usually seen in people with normal immune function, granulomatous amoebic encephalitis is usually seen in people with poor immune function, such as those with HIV/AIDS or leukemia.
Naegleriasis was the topic in Season 2 of the medical mystery drama House, M.D. in the two-part episode titled "Euphoria". It is also the topic of the episode "39 Differences" of season 6 of The Good Doctor.
Research
The U.S. National Institutes of Health budgeted $800,000 for research on the disease in 2016. Phenothiazines have been tested in vitro and in animal models of PAM. Improving case detection through increased awareness, reporting, and information about cases might enable earlier detection of infections, provide insight into the human or environmental determinants of infection, and allow improved assessment of treatment effectiveness.
| Biology and health sciences | Protozoan infections | Health |
522868 | https://en.wikipedia.org/wiki/Centaurus%20A | Centaurus A | Centaurus A (also known as NGC 5128 or Caldwell 77) is a galaxy in the constellation of Centaurus. It was discovered in 1826 by Scottish astronomer James Dunlop from his home in Parramatta, in New South Wales, Australia. There is considerable debate in the literature regarding the galaxy's fundamental properties such as its Hubble type (lenticular galaxy or a giant elliptical galaxy) and distance (11–13 million light-years). It is the closest radio galaxy to Earth, as well as the closest BL Lac object, so its active galactic nucleus has been extensively studied by professional astronomers. The galaxy is also the fifth-brightest in the sky, making it an ideal amateur astronomy target. It is only visible from the southern hemisphere and low northern latitudes.
The center of the galaxy contains a supermassive black hole with a mass of 55 million solar masses, which ejects a relativistic jet that is responsible for emissions in the X-ray and radio wavelengths. By taking radio observations of the jet separated by a decade, astronomers have determined that the inner parts of the jet are moving at about half of the speed of light. X-rays are produced farther out as the jet collides with surrounding gases, resulting in the creation of highly energetic particles. The X-ray jets of Centaurus A are thousands of light-years long, while the radio jets are over a million light-years long.
It is also one of the nearest large starburst galaxies, of which a galactic collision is suspected to be responsible for an intense burst of star formation. Models have suggested that Centaurus A was a large elliptical galaxy that collided with a smaller spiral galaxy, with which it will eventually merge. For that reason, the galaxy has been of particular interest to astronomers for years. While collisions of spiral galaxies are relatively common, the effects of a collision between an elliptical and a spiral galaxy are not fully understood.
Observational history
NGC 5128 was discovered on 29 April 1826 by James Dunlop during a survey at the Parramatta Observatory.
In 1847 John Herschel described the galaxy as "two semi-ovals of elliptically formed nebula appearing to be cut asunder and separated by a broad obscure band parallel to the larger axis of the nebula, in the midst of which a faint streak of light parallel to the sides of the cut appears."
In 1949 John Gatenby Bolton, Bruce Slee and Gordon Stanley localized NGC 5128 as one of the first extragalactic radio sources. Five years later, Walter Baade and Rudolph Minkowski suggested that the peculiar structure is the result of a merge event of a giant elliptical galaxy and a small spiral galaxy. The first detection of X-ray emissions, using a sounding rocket, was performed in 1970. In 1975–76 gamma-ray emissions from Centaurus A were observed through the atmospheric Cherenkov technique.
The Einstein Observatory detected an X-ray jet emanating from the nucleus in 1979. Ten years later, young blue stars were found along the central dust band with the Hubble Space Telescope.
The Chandra X-ray Observatory identified in 1999 more than 200 new point sources. Another space telescope, the Spitzer Space Telescope, found a parallelogram-shaped structure of dust in near infrared images of Centaurus A in 2006.
Evidence of gamma emissions with very high energy (more than 100 GeV) was detected by the H.E.S.S-Observatorium in Namibia in 2009.
The following year, Centaurus A was identified as a source of cosmic rays of highest energies, after years of observations by Pierre Auger Observatory. In 2016 a review of data from Chandra and XMM-Newton, unusual high flares of energy were found in NGC 5128 and the galaxy NGC 4636. Jimmy Irwin of University of Alabama hypothesized the discovery as potentially a black hole in a yet unknown process or an intermediate-mass black hole.
Morphology
Centaurus A may be described as having a peculiar morphology. As seen from Earth, the galaxy looks like a lenticular or elliptical galaxy with a superimposed dust lane. The peculiarity of this galaxy was first identified in 1847 by John Herschel, and the galaxy was included in Halton Arp's Atlas of Peculiar Galaxies (published in 1966) as one of the best examples of a "disturbed" galaxy with dust absorption. The galaxy's strange morphology is generally recognized as the result of a merger between two smaller galaxies.
The bulge of this galaxy is composed mainly of evolved red stars. The dusty disk, however, has been the site of more recent star formation; over 100 star formation regions have been identified in the disk.
Novae and Supernovae
Two supernovae have been detected in Centaurus A.
The first supernova, named SN 1986G, was discovered within the dark dust lane of the galaxy by Robert Evans on 3 May 1986.
It was later identified as a Type Ia supernova,
which forms when a white dwarf's mass grows large enough to ignite carbon fusion in its center, touching off a runaway thermonuclear reaction, as may happen when a white dwarf in a binary star system strips gas away from the other star. SN 1986G was used to demonstrate that the spectra of type Ia supernovae are not all identical, and that type Ia supernovae may differ in the way that they change in brightness over time.
The second supernova, designated SN 2016adj, was discovered by Backyard Observatory Supernova Search in February 2016 and was initially classified as a Type II supernova based on its H-alpha emission line.
A subsequent classification found the spectrum best resembled the Type Ib core-collapse supernova 1999dn. (See Type Ib and Ic supernovae).
In addition to these supernovae, a luminous red nova, designated AT2020nqq (type ILRT, mag. 17.8), was discovered on 27 June 2020.
Centaurus A is close enough that classical novae can also be detected. The first confirmed nova in this galaxy was discovered by BlackGEM at magnitude 18.47 on 23 December 2024, and designated AT 2024aeql.
Distance
Distance estimates to Centaurus A established since the 1980s typically range between 3–5 Mpc. Classical Cepheids discovered in the heavily obscured dust lane of Centaurus A yield a distance between ~3–3.5 Mpc, depending on the nature of the extinction law adopted and other considerations. Mira variables and Type II Cepheids were also discovered in Centaurus A, the latter being rarely detected beyond the Local Group. The distance to Centaurus A established from several indicators such as Mira variables and planetary nebulae favour a more distant value of ~3.8 Mpc.
Nearby galaxies and galaxy group information
Centaurus A is at the center of one of two subgroups within the Centaurus A/M83 Group, a nearby group of galaxies. Messier 83 (the Southern Pinwheel Galaxy) is at the center of the other subgroup. These two groups are sometimes identified as one group and sometimes identified as two groups. However, the galaxies around Centaurus A and the galaxies around M83 are physically close to each other, and both subgroups appear not to be moving relative to each other. The Centaurus A/M83 Group is located in the Virgo Supercluster.
In addition to dwarf galaxies, Centaurus A, like most galaxies, has a population of globular clusters. Some objects that appear to be globular clusters are hypothesized to be the tidally stripped cores of former galaxies. The most extreme example is the object VHH81-01, whose central black hole is estimated to be around .
Observations
Radio waves
In July 2021 the Event Horizon Telescope released a resolved image of Centaurus A showing the jet coming from the black hole at its center.
Visibility
Centaurus A is located approximately 4° north of Omega Centauri (a globular cluster visible with the naked eye). Because the galaxy has a high surface brightness and relatively large angular size, it is an ideal target for amateur astronomy observations. The bright central bulge and dark dust lane are visible even in finderscopes and large binoculars, and additional structure may be seen in larger telescopes. Claims have been made that Centaurus A is visible to the naked eye under exceptionally good conditions.
Gallery
| Physical sciences | Notable galaxies | null |
522921 | https://en.wikipedia.org/wiki/Very%20Long%20Baseline%20Array | Very Long Baseline Array | The Very Long Baseline Array (VLBA) is a system of ten radio telescopes which are operated remotely from their Array Operations Center located in Socorro, New Mexico, as a part of the National Radio Astronomy Observatory (NRAO). These ten radio antennas work together as an array that forms the longest system in the world that uses very long baseline interferometry. The longest baseline available in this interferometer is about .
The construction of the VLBA began in February 1986 and it was completed in May 1993. The first astrometrical observation using all ten antennas was carried out on May 29, 1993. The total cost of building the VLBA was about $85 million. The array is funded by the National Science Foundation, and costs about $10 million a year to operate.
Each receiver in the VLBA consists of a parabolic dish antenna 25 meters (82 feet) in diameter, along with its adjacent control building. This contains the supporting electronics and machinery for the receiver, including low-noise electronics, digital computers, data storage units, and the antenna-pointing machinery. Each of the antennas is about as tall as a ten-story building when the antenna is pointed straight up, and each antenna weighs about 218 metric tons (240 short tons).
The signals from each antenna are recorded on a bank of approximately one-terabyte hard disc drives, and the information is time-stamped using atomic clocks. Once the disc drives are loaded with information, they are carried to the Pete V. Domenici Science Operations Center at the NRAO in Socorro. There, the information undergoes signal processing in a powerful set of digital computers that carry out the interferometry. These computers also make corrections for the rotation of the Earth, the slight shifts in the crust of the Earth over time, and other small measurement errors.
Observations by the VLBA
The Very Long Baseline Array usually makes radio observations at wavelengths from three millimeters to 90 centimeters, or in other words, at frequencies from 0.3 gigahertz to 96 gigahertz. Within this frequency range, the VLBA observes in eight different frequency bands that are useful for radio astronomy. The VLBA also makes observations in two narrow radio bands below one gigahertz that include spectral lines produced by bright maser emissions.
The VLBA radio telescopes are located at:
High-Sensitivity Array
The use of the VLBA can be scheduled dynamically, and its sensitivity can be improved by a factor of five by including other radio telescopes such as the Green Bank Telescope in West Virginia, the Very Large Array (VLA) in New Mexico and the Effelsberg radio telescope in Germany. These three additional sites are brought online for as much as 100 hours per four-month trimester. In this configuration, the entire array is known as the High-Sensitivity Array (HSA). The Arecibo radio telescope in Puerto Rico was also used, before it collapsed.
Baseline distance and angular resolution
Distance between each VLBA baseline (km):
The longest baseline in the array is .
| Technology | Ground-based observatories | null |
522941 | https://en.wikipedia.org/wiki/Convention%20center | Convention center | A convention center (American English; or conference centre in British English) is a large building that is designed to hold a convention, where individuals and groups gather to promote and share common interests. Convention centers typically offer sufficient floor area to accommodate several thousand attendees. Very large venues, suitable for major trade shows, are sometimes known as exhibition halls. Convention centers typically have at least one auditorium and may also contain concert halls, lecture halls, meeting rooms, and conference rooms. Some large resort area hotels include a convention center.
In Francophone countries, the term is palais des congrès (such as the Palais des Congrès de Paris) or centre des congrès (such as the Centre des congrès de Quebec).
Types
Meeting facilities with lodging: hotels that include their own convention space in addition to accommodation and other related facilities, known as convention hotels.
Meeting facilities without lodging: are convention centers that do not include accommodation; usually located adjacent to or near a hotel(s).
Other: any convention and meeting facilities designed to hold large numbers of people. Can exist alone (e.g., stadiums, arenas, parks, etc.) or within other structures (e.g., university lecture halls, museums, theaters). Usually do not include accommodation.
History
The original convention centers or halls were in castles and palaces. Originally a hall in a castle would be designed to allow a large group of lords, knights and government officials to attend important meetings with the king. A more ancient tradition would have the king or lord decide disputes among his people. These administrative actions would be done in the great hall and would exhibit the wisdom of the king as judge to the general populace.
One of the most famous convention center debacles happened in France on June 20, 1789. King Louis XVI locked a group known as the Third Estate out of the meeting hall in Versailles. This led to the revolutionary group holding their meeting in an indoor tennis court. This was the first modern democratic conference center and lead to the Tennis Court Oath and the French Revolution.
Some historic centers
19th-century exhibition halls
1850 Bingley Hall (destroyed by fire in 1984), Birmingham, England
1851 The Crystal Palace (destroyed by fire in 1936), London, England
1855 Palais de l'Industrie (dismantled in 1897), Paris, France
1873 Alexandra Palace, London, England
1876 Memorial Hall, Philadelphia, Pennsylvania
1878 Exhibition Place, Toronto
1878 La Rural, Buenos Aires, Argentina
1878 Music Hall, Cincinnati, Ohio
1879 Garden Palace (destroyed by fire in 1882), Sydney, Australia
1880 Royal Exhibition Building, Melbourne, Australia
1898 Aberdeen Pavilion, Ottawa, Ontario
1898–1903 Beurs van Berlage, Amsterdam, Netherlands
20th-century exhibition halls
1900 Grand Palais, Paris, France
1909 Festhalle, Frankfurt, Germany
1955 McCormick Place, Chicago, Illinois
1958 Centre of New Industries and Technologies, Paris, France
1959 Las Vegas Convention Center, Las Vegas, Nevada
1974 Kenyatta International Convention Centre, Nairobi, Kenya
1975 Helsinki Fair Centre, Helsinki, Finland
1976 Georgia World Congress Center, Atlanta, Georgia
1979 Internationales Congress Centrum, Berlin, Germany
1981 Moscone Center, San Francisco, California
1983 Hong Kong Convention and Exhibition Centre, Wan Chai, Hong Kong
1985 Tampere Fair Centre, Tampere, Finland
1988 Seattle Convention Center, Seattle, Washington
1989 Taipei International Convention Center, Taipei, Taiwan
1990 Colorado Convention Center, Denver, Colorado
1993 Pennsylvania Convention Center, Philadelphia, Pennsylvania
1995 Suntec Singapore Convention and Exhibition Centre, Singapore
1997 Tokyo International Forum, Tokyo, Japan
21st-century exhibition halls
2001 Bethlehem Convention Palace, Bethlehem
2003 Walter E. Washington Convention Center, Washington, D.C.
2008 BT Convention Centre, Liverpool, England
2008 Raleigh Convention Center, Raleigh, North Carolina
2008 Taipei Nangang Exhibition Center, Taipei, Taiwan
2012 Convention Center Poet Ronaldo Cunha Lima, João Pessoa, Brazil
2014 Kaohsiung Exhibition Center, Kaohsiung, Taiwan
2017 AU Convention Center, Visakhapatnam, India
2021 Rudraksha Convention Center, Varanasi, India
2021 Bangabandhu Bangladesh–China Friendship Exhibition Center, Dhaka, Bangladesh
2024 Taoyuan Convention and Exhibition Center, Taoyuan City, Taiwan
Image gallery
| Technology | Commercial buildings | null |
303437 | https://en.wikipedia.org/wiki/Tenant%20farmer | Tenant farmer | A tenant farmer is a person (farmer or farmworker) who resides on land owned by a landlord. Tenant farming is an agricultural production system in which landowners contribute their land and often a measure of operating capital and management, while tenant farmers contribute their labor along with at times varying amounts of capital and management. Depending on the contract, tenants can make payments to the owner either of a fixed portion of the product, in cash or in a combination. The rights the tenant has over the land, the form, and measures of payment vary across systems (geographically and chronologically). In some systems, the tenant could be evicted at whim (tenancy at will); in others, the landowner and tenant sign a contract for a fixed number of years (tenancy for years or indenture). In most developed countries today, at least some restrictions are placed on the rights of landlords to evict tenants under normal circumstances.
England and Wales
Historically, rural society utilised a three-tier structure of landowners (nobility, gentry, yeomanry), tenant farmers, and farmworkers. Originally, tenant farmers were known as peasants. Under Anglo-Norman law, almost all tenants were bonded to the land, and were therefore also villeins, but after the labour shortage occasioned by the Black Death in the mid-14th century, the number of free tenants substantially increased. Many tenant farmers became affluent and socially well connected, and employed a substantial number of labourers and managed more than one farm. Tenancy could be either in perpetuity or rotated by the owners. Cottiers (cottagers) held much less land.
The 17th century to the early 19th century witnessed the growth of large estates, and the opportunity for a farmer to hold land other than by tenancy was significantly reduced, with the result that by the 19th century, about 90% of agricultural land area and holdings were tenanted, although these figures declined markedly after World War II, to around 60% in 1950 and only 35% of agricultural land area in 1994. High rates of inheritance taxes in the postwar period led to the breakup or reduction of many large estates, allowing many tenants to buy their holdings at favourable prices.
The landmark 1948 Act was enacted at a time when war-time food rationing was still in force and sought to encourage long-term investment by tenants by granting them lifetime security of tenure. Under the Agriculture (Miscellaneous Provisions) Act 1976, security was extended to spouses and relatives of tenants for two successions, providing that they had been earning the majority of their income from the holding for five years. Succession rights were however withdrawn for new tenancies in 1984 and this was consolidated in the Agricultural Holdings Act 1986. These two statutes also laid down rules for the determination of rents by the arbitration process. The 1986 statute covered tenancies over agricultural land where the land was used for a trade or business and the definition of "agriculture" in section 96(1) was wide enough to include various uses that in themselves were not agricultural but were deemed so if ancillary to agriculture (e.g. woodlands). The essence of the code was to establish complex constraints on the landlord's ability to give the notice to quit while also converting fixed-term tenancies into yearly tenancies at the conclusion of the fixed term. In addition, there was a uniform rent ascertainment scheme contained in section 12.
It became difficult to obtain new tenancies as a result of landlords' reluctance to have a tenant protected by the 1986 Act and in 1995, the government of the day, with the support of industry organizations, enacted a new market-oriented code in the form of the Agricultural Tenancies Act 1995. The protection of the 1986 Act remains in respect of tenancies created prior to the existence of the 1995 Act and for those tenancies falling within section 4 of the 1995 Act. For all other tenancies granted on or after 1 September 1995, their regulation is within the 1995 Act framework.
That Act was altered with effect from 18 October 2006 by the Regulatory Reform (Agricultural Tenancies)(England and Wales) Order 2006 SI 2006/2805, which also contains changes to the 1986 Act. Tenancies granted after 18 October 2006 over agricultural land used for a trade or business will fall within the limited protection of the 1995 Act so as to enjoy (provided the term is more than two years in length or there is a yearly tenancy) a mandatory minimum twelve months written notice to quit, including in respect of fixed terms. There is for all tenancies within the scope of the Act a mandatory tenants' right to remove fixtures and buildings (section 8) together with compensation for improvements (Part III). The rent review provisions in Part II may be the subject of choice to a much greater extent than previously. Disputes under the Act are usually, by the terms of Part IV, the subject of statutory arbitration controlled by the framework of the Arbitration Act 1996.
The current regime under the 1995 Act for regulating tenancies, commonly known as Farm Business Tenancies, permits the creation of a clearly and easily terminable interest, whether by a periodic tenancy or a fixed term. In the cycle of animal husbandry and land use and improvement, the long-term effect of the Farm Business Tenancy on the landscape of Britain is not yet proven. It was predicted by landowners and other industry spokesmen that the 1995 Act would create opportunities for new tenants by allowing large areas of new lettings but this has not happened in practice as most landowners have continued to favour share farming or management agreements over formal tenancies and the majority of new lettings under the Act have been to existing farmers, often owner-occupiers taking on extra land at significantly higher rents than could be afforded by a traditional tenant.
Canada
From the nineteenth century on, tenant farming immigrants came to Canada not just from the British Isles but also the United States of America.
Ireland
Until about 1900, the majority of Ireland was held by landlords, as much as 97% in 1870, and rented out to tenant farmers who had to pay rent to landlords and taxes to the Church of Ireland and State. The majority of the people had no access to land. 1.5% of the population owned 33.7% of the island, and 50% of the country was in the hands of only 750 families. Absenteeism was common and detrimental to the country's progress. Tenants often sub-rented small plots on a yearly basis from local farmers paying for them by labour service by a system, known as conacre, most without any lease or land rights. Irish smallholders were indistinguishable from the cottiers of England.
The abuse of tenant farmers led to widespread emigration to the United States and the colonies and was a key factor within the Home Rule Movement. They also underlined a deterioration in Protestant-Catholic relationships, although there were notable elements of cooperation in reform attempts such as the Tenant Right League of the 1850s. Following the Great Famine tenant farmers were the largest class of people. Discontent led to the Land War of the 1870s onwards, the Landlord and Tenant (Ireland) Act 1870, the founding of the Land League 1879 to establish fair rents and the fixity of tenures. The movement played a key element in the unification of country and urban classes and the creation of a national identity not existing before.
The Landlord and Tenant (Ireland) Act 1870 stands out as the first attempt to resolve problems of tenants rights in Ireland and the Land Law (Ireland) Act 1881 went even further to inspire campaigners even in Wales. The Purchase of Land (Ireland) Act 1885 followed, finally the great breakthrough after the successful 1902 Land Conference, the enactment of the Land Purchase (Ireland) Act 1903 whereby the state financed tenants to completely buy out their landlords. Under the Act of 1903 and the consequential Act of 1909, the national situation was completely transformed. When in March 1920, the Irish Estate Commission reviewed the development since 1903 under these Acts, they estimated that 83 million sterling had been advanced for transferred, whilst a further were pending costing 24 million sterling. By 1914, 75% of occupiers were buying out their landlords, mostly under the two Acts. In all, under the pre-UK Land Acts over 316,000 tenants purchased their holdings amounting to out of a total of 20 million in the country.
On the formation of the Irish Free State in 1922, the Irish Land Commission was reconstituted by the Land Law (Commission) Act, 1923. The commission had acquired and supervised the transfer of up to of farmland between 1885 and 1920 where the freehold was assigned under mortgage to tenant farmers and farm workers. The focus had been on the compulsory purchase of untenanted estates so that they could be divided into smaller units for local families. In 1983, the Commission ceased acquiring land; this signified the start of the end of the commission's reform of Irish land ownership, though freehold transfers of farmland still had to be signed off by the Commission into the 1990s. The commission was dissolved in March 1999.
Japan
In Japan, landowners turned over their land to families of tenant farmers to manage. During the Meiji period, Japanese tenant farmers were traditionally cultivators rather than capitalistic or entrepreneurial venture by nature, paid in kind for their labors. Approximately 30% of land was held by tenants. Many aspects of Tokugawa feudalism continued. After WWII, the Farm Land Reform Law of 1946 banned absentee landlordism, re-distributing land and permitted tenants to buy. By the 1950s, it virtually eliminated the landlord-tenant relationship.
Scandinavia
Historically, despite Norway being practically a Danish province for almost 300 years before 1814, the countries of Denmark, Norway, and Sweden (with Finland) had differing approaches to land tenure.
Norway
A tenant farmer in Norway was known as a husmann (plural: husmenn) and were most common in the mid-19th century when they constituted around one-quarter of the country's population. Heavy demands were placed on these tenants by their landlords, the bønder or land-owning farmers. The majority of the husmann's working hours were usually taken up by work for the landlord, leaving him little time to work on his own land or better his own situation. As a result, though the husmenn were technically free to leave the land at any time, their poor economic state made them in essence "economic serfs". Failing to own their own land also made tenant farmers ineligible to vote according to the Norwegian Constitution at the time. The number of tenant farmers in the country grew during the 19th century, rising from 48,571 in 1825 to 65,060 in 1855, the latter figure representing the height of the husmann population in Norway, most of whom lived in the eastern part of the country. Given their difficult economic and social position in Norway, many Norwegian husmenn immigrated to Canada and the United States throughout the 19th century. Following the revolutions of 1848 the husmenn's cause was taken up by Marcus Thrane. Thrane fought for the husmenn's rights at home and also encouraged them to emigrate and seek better fortunes abroad. The number of husmenn began to decline in the second half of the 19th century, and by 1910 they made up less than 5% of Norwegian society.
Sweden and Finland
The term torpare/torppari (Swedish/Finnish for crofter) refers to a slightly different type of tenant farmers, less secure than the inheritable usufruct right as åbo but sometimes with contracts as long as 50 years. The lease was, depending on the landowner's good will, in practice often transferred to a son or a widow.
Their situation was usually poor but, contrary to in Denmark, they were in theory always free to leave. The croft's lease was typically paid in the form of corvée. They would work their own land as well as that of a landowning farmer (bonde), noble or other. In some aspects their situation made them easy victims of impressment. Population growth and landreforms (enskiftet) contributed to a 19th century increase of crofts but, particularly in Sweden, also to a shift from tenant farmers to farm laborers (statare) hired on yearlong contracts, paid in-kind.
The lives of torpare and statare were described by prominent Swedish and Finnish novelists and writers such as Ivar Lo-Johansson, Jan Fridegård, Väinö Linna (Under the North Star trilogy) and Moa Martinson. The Statare system was abolished in 1918 (Finland) and 1945 (Sweden), the Torpare system more gradually.
Scotland
Scotland has its own independent legal system and the legislation there differs from that of England and Wales. Neither the AHA 1986 nor the ATA 1995 applies in Scotland. The relevant legislation for Scotland is rather the Agricultural Holdings (Scotland) Act 2003 with the following amendments in The Public Services Reform (Agricultural Holdings) (Scotland) Order 2011, The Agricultural Holdings (Amendment) (Scotland) Act 2012 and The Agricultural Holdings (Scotland) 2003 Remedial Order 2014. These supersede the previous legislation in the Agricultural Holdings (Scotland) Act 1991 and the Agriculture (Scotland) Act 1948.
For Scotland see Crofting, a traditional and long-established means of tenant and subsistence farming.
United States
Tenant farming has been important in the US from the 1870s to the present. Tenants typically bring their own tools and animals. To that extent it is distinguished from being a sharecropper, which is a tenant farmer who usually provides no capital and pays fees with crops.
A hired hand is an agricultural employee even though he or she may live on the premises and exercise a considerable amount of control over the agricultural work, such as a foreman. A sharecropper is a farm tenant who pays rent with a portion (often half) of the crop he raises and who brings little to the operation besides his family labor; the landlord usually furnishing working stock, tools, fertilizer, housing, fuel, and seed, and often providing regular advice and oversight.
Tenant farming in the North was historically a step on the "agricultural ladder" from hired hand or sharecropper taken by young farmers as they accumulated enough experience and capital to buy land (or buy out their siblings when a farm was inherited).
About two-thirds of sharecroppers were white, the rest black. Sharecroppers, the poorest of the poor, organized for better conditions. The racially integrated Southern Tenant Farmers Union made gains for sharecroppers in the 1930s. Sharecropping had diminished in the 1940s due to the Great Depression, farm mechanization, and other factors.
Black Belt conditions
In the Black Belt in the American South until the mid 20th century, the predominant agricultural system involved white land owners and African-American tenant farmers. Very little cash changed hands. The few local banks were small and cash was scarce and had to be hoarded for taxes. Landowners needed a great deal of labor at harvest time to pick the cash crop, cotton. The typical plan was to divide old plantations into small farms that were assigned to the tenants. Throughout the year the tenants lived rent-free. They tended their own gardens. Every week, they bought food and supplies on credit through the local country store. At harvest time, the tenants picked the cotton, and turned it all over to the landowners. They sold the cotton on the national market and used part of the funds to pay the debts owed to the country store. The cycle then started all over again. Landowners also worked some of the land directly, using black labor paid in cash. The landowners held all the political power, and fought vigorously against government welfare programs that would provide cash that would undermine the cashless system. Economic historians Lee Alston and Joseph Ferrie (1999) describe the system as essentially an informal contract that:
bound employer and worker through the provision of housing, medical care, and other in-kind services along with cash wages. At its heart, it guaranteed a stable and adequate labor supply to the planter. Though restricted by the directives of the planter, workers in return received some measure of economic stability, including a social safety net, access to financial capital, and physical protection in an often-violent society.
Tenant farmers often had agricultural managers who supervised their activities. In 1907, for instance, J. H. Netterville began employment for the Panola Company, an agricultural business founded by William Mackenzie Davidson in the rich farming area of St. Joseph in Tensas Parish in northeastern Louisiana in the Mississippi River delta country. In its heyday, Panola controlled some eleven thousand acres, two-thirds planted in cotton and the other third in grains. Netterville became general manager of three highly profitable Panola properties, the Balmoral, Blackwater, and Wyoming plantations near Newellton, in which capacity he supervised 125 African-American tenant farming families, with little strife and great ease, according to reports from that period.
Latin America
For tenant farmers and other landholding arrangements in Latin America, see Peasant#Latin American farmers.
In media
The Emigrants, a 1971 Swedish film in which tenant farmers emigrate to the United States
Plain Folk of the Old South, 1949 book, about U.S. before 1860
| Technology | Agriculture_2 | null |
303469 | https://en.wikipedia.org/wiki/Teak | Teak | Teak (Tectona grandis) is a tropical hardwood tree species in the family Lamiaceae. It is a large, deciduous tree that occurs in mixed hardwood forests. Tectona grandis has small, fragrant white flowers arranged in dense clusters (panicles) at the end of the branches. These flowers contain both types of reproductive organs (perfect flowers). The large, papery leaves of teak trees are often hairy on the lower surface. Teak wood has a leather-like smell when it is freshly milled and is particularly valued for its durability and water resistance. The wood is used for boat building, exterior construction, veneer, furniture, carving, turnings, and various small projects.
Tectona grandis is native to south and southeast Asia, mainly Bangladesh, India, Indonesia, Malaysia, Myanmar, Thailand, and Sri Lanka, but is naturalised and cultivated in many countries in Africa and the Caribbean. Myanmar's teak forests account for nearly half of the world's naturally occurring teak. Molecular studies show that there are two centres of the genetic origin of teak: one in India and the other in Myanmar and Laos.
Description
Teak is a large deciduous tree up to tall with grey to greyish-brown branches, known for its high-quality wood. Its leaves are ovate-elliptic to ovate, long by wide, and are held on robust petioles which are long. Leaf margins are entire.
Fragrant white flowers are borne on long by wide panicles from June to August. The corolla tube is 2.5–3 mm long with 2 mm wide obtuse lobes. Tectona grandis sets fruit from September to December; fruits are globose and 1.2–1.8 cm in diameter. Flowers are weakly protandrous in that the anthers precede the stigma in maturity and pollen is shed within a few hours of the flower opening. The flowers are primarily entomophilous (insect-pollinated), but can occasionally be anemophilous (wind-pollinated). A 1996 study found that in its native range in Thailand, the major pollinators were species in the bee genus Ceratina.
Wood
Wood texture is hard and rings are porous. The density varies according to moisture content: at 15% moisture content it is 660 kg/m. The heartwood is yellowish to golden-brown. Sapwood is whitish to pale yellowish brown. It can easily separate from heartwood. Teak darkens as it ages. There can be a large variation, depending on which region the teak is from. Old growth has much tighter rings than new growth. There is a leather-like scent in newly cut wood.
Botanical history
Tectona grandis was first formally described by Carl Linnaeus the Younger in his 1782 work Supplementum Plantarum. In 1975, Harold Norman Moldenke published new descriptions of four forms of this species in the journal Phytologia. Moldenke described each form as varying slightly from the type specimen: T. grandis f. canescens is distinguished from the type material by being densely canescent or covered in hairs, on the underside of the leaf, T. grandis f. pilosula is distinct from the type material in the varying morphology of the leaf veins, T. grandis f. punctata is only hairy on the larger veins on the underside of the leaf, and T. grandis f. tomentella is noted for its dense yellowish tomentose hairs on the lower surface of the leaf.
Etymology
The English word teak comes via the Portuguese from Malayalam (cognate with Tamil , Telugu , and Kannada ) via Sanskrit "shaka" and "saka". Central Province teak and Nagpur teak are named for those regions of India.
Distribution and habitat
Tectona grandis is one of three species in the genus Tectona. The other two species, T. hamiltoniana and T. philippinensis, are endemics with relatively small native distributions in Myanmar and the Philippines, respectively. Tectona grandis is native to India, Bangladesh, Sri Lanka, Indonesia, Myanmar, northern Thailand, and northwestern Laos.
Tectona grandis is found in a variety of habitats and climatic conditions from arid areas with only 500 mm of rain per year to very moist forests with up to 5,000 mm of rain per year. Typically, though, the annual rainfall in areas where teak grows averages 1,250–1,650 mm with a 3–5 month dry season.
Cultivation
Teak's natural oils make it useful in exposed locations and make the timber termite- and pest-resistant. Teak is durable even when not treated with oil or varnish. Timber cut from old teak trees was once believed to be more durable and harder than plantation-grown teak. Studies have shown that plantation teak performs on par with old-growth teak in erosion rate, dimensional stability, warping, and surface checking, but is more susceptible to colour change from UV exposure.
The vast majority of commercially harvested teak is grown on teak plantations found in Indonesia and controlled by Perum Perhutani (a state-owned forest enterprise) that manages the country's forests. The primary use of teak harvested in Indonesia is in the production of outdoor teak furniture for export. Nilambur in Kerala, India, is also a major producer of teak and is home to the world's oldest teak plantation.
Teak consumption raises several environmental concerns, such as the disappearance of rare old-growth teak. However, its popularity has led to growth in sustainable plantation teak production throughout the seasonally dry tropics in forestry plantations. The Forest Stewardship Council offers certification of sustainably grown and harvested teak products. Propagation of teak via tissue culture for plantation purposes is commercially viable.
Teak plantations were widely established in Equatorial Africa during the Colonial era. These timber resources, as well as the oil reserves, are at the heart of the current (2014) South Sudanese conflict.
Much of the world's teak is exported by Indonesia and Myanmar. There is also a rapidly growing plantation-grown market in Central America (Costa Rica) and South America. With a depletion of remaining natural hectares of teak forests, growth in plantations in Latin America is expected to rise.
Hyblaea puera, commonly known as the teak defoliator, is a moth native to southeast Asia. It is a teak pest whose caterpillar feeds on teak and other species of trees common in the region of Southeast Asia.
Uses
Teak's high oil content, high tensile strength, and tight grain make it particularly suitable where weather resistance is desired. It is used in the manufacture of outdoor furniture and boat decks. It is also used for cutting boards, indoor flooring, countertops, and as a veneer for indoor finishings. Although easily worked, it can cause severe blunting on edged tools because of the presence of silica in the wood. Over time teak can weather to a silvery-grey finish, especially when exposed to sunlight.
Teak is used extensively in India to make doors and window frames, furniture, and columns, and beams in homes. It is resistant to termite attacks and damage caused by other insects. Mature teak fetches a very good price. It is grown extensively by forest departments of different states in forest areas. It was also used in the construction of the Kaaba in the Masjid al-Haram of Mecca, which is the holiest structure in the Islamic faith.
Leaves of the teak wood tree are used in making Pellakai gatti (jackfruit dumpling), where batter is poured into a teak leaf and steamed. This type of usage is found in the coastal district of Udupi in the Tulunadu region in South India. The leaves are also used in gudeg, a dish of young jackfruit made in Central Java, Indonesia, and give the dish its dark brown colour.
Teak is used as a food plant by the larvae of moths of the genus Endoclita including E. aroura, E. chalybeatus, E. damor, E. gmelina, E. malabaricus, E. sericeus and E. signifer other Lepidoptera including the turnip moth.
Boatbuilding
Teak has been used as a boatbuilding material for over 2000 years (it was found in an archaeological dig in Berenice Panchrysos, a port on the Indian Roman trade route). In addition to relatively high strength, teak is also highly resistant to rot, fungi, and mildew. The wood has a relatively low shrinkage ratio, which makes it excellent for applications where it undergoes periodic changes in moisture. Teak has the unusual property of being both an excellent structural timber for framing or planking, while at the same time being easily worked and finished, unlike some otherwise similar woods such as purpleheart. For this reason, it is also prized for the trim work on boat interiors. Due to the oily nature of the wood, care must be taken to properly prepare the wood before gluing.
When used on boats, teak is also very flexible in the finishes that may be applied. One option is to use no finish at all, in which case the wood will naturally weather to a pleasing silver grey. The wood may also be oiled with a finishing agent such as linseed or tung oil. This results in a somewhat dull finish. Finally, teak may also be varnished for a deep, lustrous glow.
Teak is also used extensively in boat decks, as it is extremely durable but requires regular maintenance. The teak tends to wear into the softer 'summer' growth bands first, forming a natural 'non-slip' surface. Any sanding is therefore only damaging. The use of modern cleaning compounds, oils or preservatives will shorten the life of the teak, as it contains natural teak oil a very small distance below the white surface. Wooden boat experts will only wash the teak with salt water, and re-caulk when needed. This cleans the deck and prevents it from drying out and the wood shrinking. The salt helps it absorb and retain moisture and prevents any mildew and algal growth. Over-maintenance, such as cleaning teak with harsh chemicals, can shorten its usable lifespan as decking.
Propagation
Teak is propagated mainly from seeds. Germination of the seeds involves pretreatment to remove dormancy arising from the thick pericarp. Pretreatment involves alternate wetting and drying of the seed. The seeds are soaked in water for 12 hours and then spread to dry in the sun for 12 hours. This is repeated for 10–14 days and then the seeds are sown in shallow germination beds of coarse peat covered by sand. The seeds then germinate after 15 to 30 days.
Clonal propagation of teak has been successfully done through grafting, rooted stem cuttings, and micropropagation. While bud grafting onto seedling root stock has been the method used for establishing clonal seed orchards that enables assemblage of clones of the superior trees to encourage crossing, rooted stem cuttings and micro propagated plants are being increasingly used around the world for raising clonal plantations.
Illegal logging
Illegal logging is prevalent in countries with natural teak forests, including India and Burma.
Since 1989, the state-owned Myanma Timber Enterprise has run the country's logging industry. In 2014, Myanmar's government imposed a strict ban on exporting wild-grown teak logs. In 2015, 153 Chinese loggers were sentenced to life in prison for illegal logging. Illegal teak logging persists, especially in contested areas. While it is illegal for timber to be exported via land borders, 95% of Myanmar's teak enters China through the China–Myanmar border.
Since the 2021 Myanmar coup d'état, illegal logging of teak and tamalan trees has surged in Sagaing Region, predominantly in key contested battlegrounds, including Kani, Yinmabin, Kantbalu, Indaw, and Banmauk townships. Both the Burmese military and resistance groups have profited from the illegal logging trade. Smugglers transport the wood to India to circumvent economic sanctions and use the Myanma Timber Enterprise to license the wood as being sourced from permitted areas.
EU regulation
The regulation that addresses the import of timber, including teak, into the EU from unknown or illegal sources is the EU Timber Regulation (EUTR) No. 995/2010. This regulation aims to prevent the trade of illegally harvested timber and timber products within the EU market. It places an obligation on operators who place timber and timber products on the EU market to ensure they are legally harvested.
This regulation specifically applies to teak and other high-risk timber species, particularly those sourced from countries with poor forest governance or illegal logging activities. Myanmar, for example, has been a focus due to concerns over illegal teak harvesting from there.
World's largest living teak tree
Ministry of Environmental Conservation and Forestry (Myanmar) found the world's two biggest living teak trees on 28 August 2017 in Homalin Township, Sagaing Region, Myanmar. The biggest one, named Homemalynn 1, is in diameter and tall. The second biggest one, named Homemalynn 2, is in diameter.
Previously, the world's biggest recorded teak tree was located within the Parambikulam Wildlife Sanctuary in the Palakkad District of Kerala in India, named Kannimara. The tree is approximately tall. Its age is between 450 and 500 years and is considered one of the oldest teak trees in the world.
In 2017, a tree was discovered in the Ottakallan area of the Thundathil range of the Malayattoor Forest Division in Kerala with a girth of and a height of . A teak tree in Kappayam, Edamalayar, Kerala, which used to be considered the biggest, has a girth of 7.23 metres.
Tree No. 23 is the oldest planted teak on earth. It is located in Conolly's plot (the world's oldest teak plantation), Nilambur, Kerala.
| Biology and health sciences | Lamiales | Plants |
303612 | https://en.wikipedia.org/wiki/Homo%20rudolfensis | Homo rudolfensis | Homo rudolfensis is an extinct species of archaic human from the Early Pleistocene of East Africa about 2 million years ago (mya). Because H. rudolfensis coexisted with several other hominins, it is debated what specimens can be confidently assigned to this species beyond the lectotype skull KNM-ER 1470 and other partial skull aspects. No bodily remains are definitively assigned to H. rudolfensis. Consequently, both its generic classification and validity are debated without any wide consensus, with some recommending the species to actually belong to the genus Australopithecus as A. rudolfensis or Kenyanthropus as K. rudolfensis, or that it is synonymous with the contemporaneous and anatomically similar H. habilis.
H. rudolfensis is distinguished from H. habilis by larger size, but it is also argued that this species actually consists of male H. habilis specimens, assuming that H. habilis was sexually dimorphic and males were much larger than females. Because no bodily remains are definitely identified, body size estimates are largely based on the stature of H. habilis. Using this, male H. rudolfensis may have averaged about in height and in weight, and females and . KNM-ER 1470 had a brain volume of about . Like other early Homo, H. rudolfensis had large cheek teeth and thick enamel.
Early Homo species exhibit marked brain growth compared to Australopithecus predecessors, which is typically explained as a change in diet with a calorie-rich food source, namely meat. Though not associated with tools, dental anatomy suggests some processing of plant or meat fiber before consumption, though the mouth could still effectively chew through mechanically challenging food, indicating tool use did not greatly affect diet.
Research history
The first fossils were discovered in 1972 along Lake Turkana (at the time called Lake Rudolf) in Kenya, and were detailed by Kenyan palaeoanthropologist Richard Leakey the following year. The specimens were: a large and nearly complete skull (KNM-ER 1470, the lectotype) discovered by Bernard Ngeneo, a local; a right femur (KNM-ER 1472) discovered by J. Harris; an upper femur (proximal) fragment (KNM-ER 1475) discovered by fossil collector Kamoya Kimeu; and a complete left femur (KNM-ER 1481) discovered by Harris. However, it is unclear if the femora belong to the same species as the skull. Leakey classified them under the genus Homo because he had reconstructed the skull fragments so that it had a large brain volume and a flat face, but did not assign them to a species. Because the horizon they were discovered in was, at the time, dated to 2.9–2.6 million years ago (mya), Leakey thought these specimens were a very early human ancestor. This challenged the major model of human evolution at the time where Australopithecus africanus gave rise to Homo about 2.5 mya, but if Homo had already existed at this time, it would call for serious revisions. However, the area was redated to about 2 mya in 1977 (the same time period as H. habilis and H. ergaster/H. erectus), and more precisely to 2.1–1.95 mya in 2012. They were first assigned to the species habilis in 1975 by anthropologists Colin Groves and Vratislav Mazák. In 1978, in a joint paper with Leakey and English anthropologist Alan Walker, Walker suggested the remains belong in Australopithecus (and that the skull was incorrectly reconstructed), but Leakey still believed they belonged to Homo, though they both agreed that the remains could belong to habilis.
KNM-ER 1470 was much larger than the Olduvai remains, so the terms H. habilis sensu lato ("in the broad sense") and H. habilis sensu stricto ("in the strict sense") were used to include or exclude the larger morph, respectively. In 1986, English palaeoanthropologist Bernard Wood first suggested these remains represent a different Homo species, which coexisted with H. habilis and H. ergaster/H. erectus. Coexisting Homo species conflicted with the predominant model of human evolution at the time which was that modern humans evolved in a straight line directly from H. ergaster/H. erectus which evolved directly from H. habilis. In 1986, the remains were placed into a new species, rudolfensis, by Russian anthropologist Valery Alekseyev (but he used the genus Pithecanthropus, which was changed to Homo three years later by Groves). In 1999, Kennedy argued that the name was invalid because Alekseyev had not assigned a holotype. Pointing out that this is in fact not mandatory, Wood the same year nevertheless designated KNM-ER 1470 as the lectotype. However, the validity of this species has also been debated on material grounds, with some arguing that H. habilis was highly sexually dimorphic like modern non-human apes, with the larger skulls classified as "H. rudolfensis" actually representing male H. habilis. In 1999, Wood and biological anthropologist Mark Collard recommended moving rudolfensis and habilis to Australopithecus based on the similarity of dental adaptations. However, they conceded that dental anatomy is highly variable among hominins and not always reliable when formulating family trees.
In 2003, Australian anthropologist David Cameron concluded that the earlier australopithecine Kenyanthropus platyops was the ancestor of rudolfensis, and reclassified it as K. rudolfensis. He also believed that Kenyanthropus was more closely related to Paranthropus than Homo. In 2008, a re-reconstruction of the skull concluded it was incorrectly restored originally, though agreed with the classification as H. rudolfensis. In 2012, British palaeoanthropologist Meave Leakey described the juvenile partial face KNM-ER 62000 discovered in Koobi Fora, Kenya; noting it shares several similarities to KNM-ER 1470 and is smaller, she assigned it to H. rudolfensis, and, because prepubescent male and female bones should be indistinguishable, differences between juvenile H. rudolfensis and adult H. habilis specimens support species distinction. She also concluded that the jawbone KNM-ER 1802, an important specimen often used in classifying other specimens as H. rudolfensis, actually belongs to a different (possibly undescribed) species, but American palaeoanthropologist Tim D. White believes this to be premature because it is unclear how wide the range of variation is in early hominins. The 2013 discovery of the 1.8 Ma Georgian Dmanisi skulls which exhibit several similarities with early Homo have led to suggestions that all contemporary groups of early Homo in Africa, including H. habilis and H. rudolfensis, are the same species and should be assigned to H. erectus. There is still no wide consensus on how rudolfensis and habilis relate to H. ergaster and descendent species.
Beyond KNM-ER 1470, there is disagreement on which specimens actually belong in H. rudolfensis as it is difficult to assign with accuracy remains that do not preserve the face and jaw. No H. rudolfensis bodily elements have been definitively associated with a skull and thus to the species. Most proposed H. rudolfensis fossils come from Koobi Fora and date to 1.9–1.85 mya. Remains from the Shungura Formation, Ethiopia, and Uraha, Malawi, are dated as far back as 2.5–2.4 mya, which would make it the earliest identified species of Homo. The latest potential specimen is KNM-ER 819 dating to 1.65–1.55 mya.
Nonetheless, H. rudolfensis and H. habilis generally are recognised members of the genus at the base of the family tree, with arguments for synonymisation or removal from the genus not widely adopted. Though it is now largely agreed upon that Homo evolved from Australopithecus, the timing and placement of this split has been much debated, with many Australopithecus species having been proposed as the ancestor. The discovery of LD 350-1, the oldest Homo specimen, dating to 2.8 mya, in the Afar Region of Ethiopia may indicate that the genus evolved from A. afarensis around this time. The species LD 350-1 belongs to could be the ancestor of H. rudolfensis and H. habilis, but this is unclear. Based on 2.1 million year old stone tools from Shangchen, China, possibly an ancestral species to H. rudolfensis and H. habilis dispersed across Asia.
Anatomy
Skull
In 1973, Richard Leakey had reconstructed the skull KNM-ER 1470 with a flat face and a brain volume of . In 1983, American physical anthropologist Ralph Holloway revised the base of the skull and calculated a brain volume of . For comparison, H. habilis specimens average about , and H. ergaster . Anthropologist Timothy Bromage and colleagues revised the face again at a 5° incline (slightly prognathic) instead of completely flat, but pushed the nasal bone back directly beneath the frontal bones. He then said it was possible to predict brain size based on just the face and (disregarding the braincase) calculated , and chalked up the errors of Leakey's reconstruction to a lack of research of the biological principles of facial anatomy at the time as well as confirmation bias, as a flat-faced reconstruction of the skull aligned with the predominant model of human evolution at the time. This was refuted by American palaeoanthropologist John D. Hawks because the skull remained more or less unchanged except for the 5° rotation outwards. Bromage and colleagues returned in 2008 with a revised skull reconstruction and brain volume estimate of .
Fossils have generally been classified into H. rudolfensis due to large skull size, flatter and broader face, broader cheek teeth, more complex tooth crowns and roots, and thicker enamel compared to H. habilis. Early Homo are characterised by larger teeth compared to later Homo. The cheek teeth of KNM-ER 60000, a jawbone, in terms of size are on the lower end for early Homo, except for the third molar which is within range. The molars increase in size towards the back of the mouth. The tooth rows of KNM-ER 1470, KNM-ER 60000, and KNM-ER 62000 are rectangular, whereas the tooth row of KNM-ER 1802 is U-shaped, which may indicate that these two morphs represent different species, or demonstrate the normal range of variation for H. rudolfensis jaws. In UR 501 from Uraha, Malawi—the oldest H. rudolfensis specimen dating to 2.5–2.3 mya—the tooth enamel thickness is the same as in other early Homo, but the enamel on the molars is almost as thick as Paranthropus molars (which have some of the thickest enamel of any hominin). Such a wide variation in enamel thickness across the cheek teeth is not exhibited in KNM-ER 1802, which may indicate regional differences among H. rudolfensis populations.
Build
Body size estimates of H. rudolfensis and H. habilis typically conclude a small size comparable to australopithecines. These largely depend on the H. habilis partial skeleton OH 62 estimated at in height and in weight. H. rudolfensis is thought to be bigger than H. habilis, but it is unclear how big this species was as no bodily elements have been definitively associated with a skull. Based on just the KNM-ER 1470 skull, male H. rudolfensis were estimated to have been in height and in weight, and females and .
For specimens that might be H. rudolfensis: the femur KNM-ER 1472 which may also be H. habilis or H. ergaster was estimated at and , the humerus KNM-ER 1473 and , the partial leg KNM-ER 1481 which may also be H. ergaster and , the pelvis KNM-ER 3228 which may also be H. ergaster and , and the femur KNM-ER 3728 which may be H. habilis or P. boisei and . It is generally assumed that pre-H. ergaster hominins, including H. rudolfensis and H. habilis, exhibited sexual dimorphism with males markedly bigger than females. However, relative female body mass is unknown in either species.
Early hominins, including H. rudolfensis, are thought to have had thick body hair coverage like modern non-human apes because they appear to have inhabited cooler regions and are thought to have had a less active lifestyle than (presumed hairless) post-ergaster species, and so probably required thick body hair to stay warm. The juvenile specimen KNM-ER 62000, a partial face, has the same age landmarks as a 13 to 14 year old modern human, but more likely died at around 8 years of age due to the presumed faster growth rate among early hominins based on dental development rate.
Culture
It is typically thought that the diets of early Homo had a greater proportion of meat than Australopithecus, and that this led to brain growth. The main hypotheses regarding this are: meat is energy- and nutrient-rich and put evolutionary pressure on developing enhanced cognitive skills to facilitate strategic scavenging and monopolise fresh carcasses, or meat allowed the large and calorie-expensive ape gut to decrease in size allowing this energy to be diverted to brain growth. Alternatively, it is also suggested that early Homo, in a drying climate with scarcer food options, relied primarily on underground storage organs (such as tubers) and food sharing, which facilitated social bonding among both male and female group members. However, unlike what is presumed for H. ergaster and later Homo, short-statured early Homo were likely incapable of endurance running and hunting, and the long and Australopithecus-like forearm of H. habilis could indicate early Homo were still arboreal to a degree. Also, organised hunting and gathering is thought to have emerged in H. ergaster. Nonetheless, the proposed food-gathering models to explain large brain growth necessitate increased daily travel distance. Large incisor size in H. rudolfensis and H. habilis compared to Australopithecus predecessors implies these two species relied on incisors more. The large, Australopithecus-like molars could indicate more mechanically challenging food compared to later Homo. The bodies of the mandibles of H. rudolfensis and other early Homo are thicker than those of modern humans and all living apes, more comparable to Australopithecus. The mandibular body resists torsion from the bite force or chewing, meaning their jaws could produce unusually powerful stresses while eating.
H. rudolfensis is not associated with any tools. However, the greater molar cusp relief in H. rudolfensis and H. habilis compared to Australopithecus suggests the former two used tools to fracture tough foods (such as pliable plant parts or meat), otherwise the cusps would have been more worn down. Nonetheless, the jaw adaptations for processing mechanically challenging food indicates technological advancement did not greatly affect their diet. Large concentrations of stone tools are known from Koobi Fora. Because these aggregations are coincident with the emergence of H. ergaster, it is probable H. ergaster manufactured them, though it is not possible to definitively attribute the tools to a species because H. rudolfensis, H. habilis, and P. boisei are also well known from the area. The oldest specimen of Homo, LD 350-1, is associated with the Oldowan stone tool industry, meaning this tradition had been in use by the genus since near its emergence.
Early H. rudolfensis and Paranthropus have exceptionally thick molars for hominins, and the emergence of these two coincides with a cooling and aridity trend in Africa about 2.5 mya. This could mean they evolved due to climate change. Nonetheless, in East Africa, tropical forests and woodlands still persisted through periods of drought. H. rudolfensis coexisted with H. habilis, H. ergaster, and P. boisei.
| Biology and health sciences | Homo | Biology |
303731 | https://en.wikipedia.org/wiki/Risso%27s%20dolphin | Risso's dolphin | Risso's dolphin (Grampus griseus) is a marine mammal and dolphin, the only species of the genus Grampus. Some of the most closely related species to these dolphins include: pilot whales (Globicephala spp.), pygmy killer whales (Feresa attenuata), melon-headed whales (Peponocephala electra), and false killer whales (Pseudorca crassidens). These dolphins grow to be about 10 ft in length and can be identified by heavy scarring that appears white. They are located worldwide in cold to temperate waters, but most typically found along continental shelves due to their eating habits. Risso's dolphins have a diet that contains primarily cephalopods. They are able to search for prey at various depths due to their ability to reach depths of almost 600m. Individuals typically travel in pods ranging anywhere from 10 to 50 dolphins, with which they form tight social bonds.
Along with most marine species, the Risso's dolphin suffers from anthropogenic disruptions to the environment. Pollution, both from noise and plastics, is a common cause of higher mortality rates. Many can be, or have been, affected by entanglement in fishing nets and whaling. Risso's dolphins are currently protected in the United States; however, they are still hunted in other parts of the world.
Taxonomy
Risso's dolphin is named after Antoine Risso, whose study of the animal formed the basis of the recognized description by Georges Cuvier in 1812. The holotype referred to specimen at the Muséum National d'Histoire Naturelle, an exhibit using preserved skin and skull obtained at Brest, France.
The type and sole species of the genus Grampus refers to Delphinus griseus Cuvier 1812. A proposition to name this genus Grampidelphis in 1933, when the taxonomic status of 'blackfish' was uncertain, and conserving the extensive use of "Grampus" for the 'killer' Orcinus orca", also suggested renaming this species (Grampidelphis exilis Iredale, T. & Troughton, E. le G. 1933). These were recognised as synonyms after publication of the Catalog of Whales (Hershkovitz, 1966).
Another common name for the Risso's dolphin is grampus (also the species' genus), although this common name was more often used for the orca. The etymology of the word "grampus" is unclear. It may be an agglomeration of the Latin or French , both meaning big fish. The specific epithet griseus refers to the mottled (almost scarred) grey colour of its body.
Description
Risso's dolphin has a relatively large anterior body and dorsal fin, while the posterior tapers to a relatively narrow tail. The bulbous shape of the head has a vertical crease in front.
Infants are dorsally grey to brown and ventrally cream-colored, with a white anchor-shaped area between the pectorals and around the mouth. In older calves, the nonwhite areas darken to nearly black, and then lighten (except for the always dark dorsal fin). Linear scars mostly from social interaction eventually cover the bulk of the body; scarring is a common feature of male to male competition in toothed whales, but Risso's dolphin tend to be unusually heavily scarred. The pronounced appearance of these scars results from the lack of repigmentation, which may be advantageous as a display that reduces further challenges from other males. Older individuals appear mostly white. Most individuals have two to seven pairs of teeth, all in the lower jaw.
Length is typically , although specimens may reach . Like most dolphins, males are typically slightly larger than females. This species weighs , making it the largest species called "dolphin".
Range and habitat
Risso's dolphins are found nearly worldwide, from cold and temperate to tropical waters, in the Indian, Pacific and Atlantic Oceans, as well as parts of the Baltic Sea, the Persian Gulf and the Mediterranean, North and Red Seas (excepting the Black Sea; however, a rare stranding was recorded in the Sea of Marmara in 2012).). There have been several documented sightings in Roskilde Fjord, in the waters of Lejre Vig, just off of the coast of Skjoldungernes Land National Park, Denmark. Analysis of Risso's dolphins found in the U.K. and in the Mediterranean display variations in mitochondrial DNA. It is possible that one reason for these differences could be the lack of interaction between individuals in the two locations.
In the Pacific, they range from French Polynesia west to Samoa, north to the Hawaiian Islands, as far as the Gulf of Alaska. However, they are absent from the waters of the western Pacific (off of Asia) beyond Futuna. They are quite common along the western coasts of British Columbia, the United States and Mexico, continuing their range to the southern tip of Tierra Del Fuego. In the eastern Atlantic, they have been sighted as far south as the offshore waters of Liberia, Guinea and Western Africa north through the Canary Islands and the Azores to southern Greenland. On the western Atlantic side, Risso’s dolphins have been seen as far south as Guyana and Martinique; they can be found throughout much of the Caribbean Sea and the Gulf of Mexico to Florida and the Bahamas, and all along the American Eastern Seaboard and the Canadian Maritime Provinces.
Their preferred environment is just off the continental shelf, on steep banks, with water depths varying from , and water temperatures at least and preferably . They have been recorded diving to depths of up to in pursuit of prey.
Since at least 2017, Risso's dolphins have begun to appear off of the subarctic Norwegian coast, as far north as Bleik's Canyon, off of Andøya. The repeated, regular sightings imply an expansion of their natural range. Possible explanations for this movement are a changing climate or varying water currents, as well as a northward migration of prey species or competition with other cetaceans, such as pilot whales.
Due to the low population density of the species, Risso's dolphins are widely considered difficult to establish an accurate estimate of population size in any given area.
Ecology
They feed almost exclusively on neritic and oceanic squid, mostly nocturnally. Predation does not appear significant. Mass strandings are infrequent. Analysis carried out on the stomach contents of stranded specimens in Scotland showed that the most important species preyed on in Scottish waters is the curled octopus (Eledone cirrhosa).
A population is found off Santa Catalina Island where they are sympatric with short-finned pilot whales (Globicephala macrorhynchus) and both species feed on the squid population. Although these species have not been seen to interact with each other, they take advantage of the commercial squid fishing that takes place at night. They have been seen by fishermen to feed around their boats. They also travel with other cetaceans. They surf the bow waves of gray whales, as well as ocean swells.
Risso's dolphins have a stratified social organisation. These dolphins typically travel in groups of between 10 and 51, but can sometimes form "super-pods" reaching up to a few thousand individuals. Smaller, stable subgroups exist within larger groups. These groups tend to be similar in age or sex. Risso's experience fidelity towards their groups. Long-term bonds are seen to correlate with adult males. Younger individuals experience less fidelity and can leave and join groups. Mothers show a high fidelity towards a group of mother and calves, but it is unclear whether or not these females stay together after their calves leave or remain in their natal pods.
Behavior
Feeding
Like many dolphin species, they use echolocation to target cephalapods and fish that are feeding below. Tagging of a population in the Azores revealed that Grampus griseus plan whether to make a shallow or deep dive, with different strategies that create profitable foraging for the considerable expenditure in time and energy. Risso's can achieve depths over by exhausting their lungs and using several spins to rapidly descend, almost vertically, and increase the time spent foraging. This allows the species to exploit a deep and dispersed layer of prey such as squid, those taking refuge during daylight when they become more vulnerable to predation. When feeding in shallow depths, however, Risso's can experience competition from other cetaceans.
Social behavior
Risso's dolphins are known to have a very active surface presence, often either displaying their tail flukes and pectoral fins, or slapping the surface of the water. They have also been known to engage in a behavior called spy-hopping, a common behavior in cetaceans where an individual vertically pokes their head out of the water. Recent studies have discussed the possibilities of spy-hopping as a sexual behavior, as it is typically only done in the presence of other individuals. Risso's dolphins do not require cutting teeth to process their cephalopod prey, which has allowed the species to evolve teeth as display weapons in mating conflicts.
Reproduction
Gestation requires an estimated 13–14 months, at intervals of 2.4 years. Calving reaches seasonal peaks in the winter in the eastern Pacific and in the summer and fall in the western Pacific. Females mature sexually at ages 8–10, and males at age 10–12. The oldest specimen reached 39.6 years.
Risso's dolphins have successfully been taken into captivity in Japan and the United States, although not with the regularity of bottlenose dolphins or orcas.
Recent studies have shown a possibility of hybridization occurrences between Risso's dolphins and bottlenose dolphins. So far, there have been for possible hybrid individuals documented in United Kingdom waters. Hybridization is not something that is uncommon with cetaceans, so it is likely that these hybrids do not have any evolutionary advantage but instead are more likely an uncommon chance event.
Human interactions
Like other dolphins and marine animals, there have been documentations of these dolphins getting caught in seine-nets and gillnets across the globe. Many of these incidents have resulted in death. Small whaling operations have also been cause of some of these deaths. Pollution has also affected many individuals who have ingested plastic. Samples from these animals shows contamination within their tissue.
Increasing oceanic noise due to human presence in the ocean threatens populations of Risso's dolphins. The intensity of anthropogenic noise can push dolphins to strand themselves and to leave their typical habitats. Evidence shows that motorized vessels create a low-frequency noise that disrupt typical acoustic behavior. This behavior is measured by the regular click trains, buzzes, pulses, and barks. The click trains produced by Risso's dolphins are necessary for Risso's dolphins to navigate through their environment and identify prey. Barks are more often used in social settings.
In Ireland, though not apparently in England, Risso's Dolphin was one of the royal fish which by virtue of the royal prerogative were the exclusive property of the English Crown.
A famed individual named Pelorus Jack was widely reported between 1888 and 1912, travelling with ships navigating the Cook Strait in New Zealand. A law protecting the animal was passed after a public outcry, renewed twice more, but suggested be invalid by its reference to Fisheries acts that did not concern marine mammals.
Conservation
The Risso's dolphin populations of the North, Baltic, and Mediterranean Seas are listed on Appendix II of the Convention on the Conservation of Migratory Species of Wild Animals (CMS), since they have an unfavourable conservation status or would benefit significantly from international co-operation organised by tailored agreements.
In addition, Risso's dolphin is covered by the Agreement on the Conservation of Small Cetaceans of the Baltic, North East Atlantic, Irish and North Seas (ASCOBANS), the Agreement on the Conservation of Cetaceans in the Black Sea, Mediterranean Sea and Contiguous Atlantic Area (ACCOBAMS), the Memorandum of Understanding for the Conservation of Cetaceans and Their Habitats in the Pacific Islands Region (Pacific Cetaceans MoU) and the Memorandum of Understanding Concerning the Conservation of the Manatee and Small Cetaceans of Western Africa and Macaronesia (Western African Aquatic Mammals MoU).
Risso's dolphins are protected in the United States under the Marine Mammal Protection Act of 1992. Currently, Japan, Indonesia, the Solomon Islands, and The Lesser Antilles hunt Risso's dolphins.
Strandings
At least one case report of strandings in Japan's Goto Islands has been associated with parasitic neuropathy of the eighth cranial nerve by a trematode in the genus Nasitrema. There was a recent reporting of a juvenile male Risso's dolphin that was stranded alive on the coast of Gran Canaria on 26 April 2019. This was the first documented case of capture myopathy and stress cardiomyopathy in a male juvenile Risso's dolphin that has received rehabilitation.
| Biology and health sciences | Toothed whale | Animals |
303802 | https://en.wikipedia.org/wiki/Phase%20rule | Phase rule | In thermodynamics, the phase rule is a general principle governing multi-component, multi-phase systems in thermodynamic equilibrium. For a system without chemical reactions, it relates the number of freely varying intensive properties () to the number of components (), the number of phases (), and number of ways of performing work on the system ():
Examples of intensive properties that count toward are the temperature and pressure. For simple liquids and gases, pressure-volume work is the only type of work, in which case .
The rule was derived by American physicist Josiah Willard Gibbs in his landmark paper titled On the Equilibrium of Heterogeneous Substances, published in parts between 1875 and 1878.
The number of degrees of freedom (also called the variance) is the number of independent intensive properties, i.e., the largest number of thermodynamic parameters such as temperature or pressure that can be varied simultaneously and independently of each other.
An example of a one-component system () is a pure chemical. A two-component system () has two chemically independent components, like a mixture of water and ethanol. Examples of phases that count toward are solids, liquids and gases.
Foundations
A phase is a form of matter that is homogeneous in chemical composition and physical state. Typical phases are solid, liquid and gas. Two immiscible liquids (or liquid mixtures with different compositions) separated by a distinct boundary are counted as two different phases, as are two immiscible solids.
The number of components (C) is the number of chemically independent constituents of the system, i.e. the minimum number of independent species necessary to define the composition of all phases of the system.
The number of degrees of freedom (F) in this context is the number of intensive variables which are independent of each other.
The basis for the rule is that equilibrium between phases places a constraint on the intensive variables. More rigorously, since the phases are in thermodynamic equilibrium with each other, the chemical potentials of the phases must be equal. The number of equality relationships determines the number of degrees of freedom. For example, if the chemical potentials of a liquid and of its vapour depend on temperature (T) and pressure (p), the equality of chemical potentials will mean that each of those variables will be dependent on the other. Mathematically, the equation , where , the chemical potential, defines temperature as a function of pressure or vice versa. (Caution: do not confuse as pressure with , number of phases.)
To be more specific, the composition of each phase is determined by intensive variables (such as mole fractions) in each phase. The total number of variables is , where the extra two are temperature T and pressure p. The number of constraints is , since the chemical potential of each component must be equal in all phases. Subtract the number of constraints from the number of variables to obtain the number of degrees of freedom as .
The rule is valid provided the equilibrium between phases is not influenced by gravitational, electrical or magnetic forces, or by surface area, and only by temperature, pressure, and concentration.
Consequences and examples
Pure substances (one component)
For pure substances so that . In a single phase () condition of a pure component system, two variables (), such as temperature and pressure, can be chosen independently to be any pair of values consistent with the phase. However, if the temperature and pressure combination ranges to a point where the pure component undergoes a separation into two phases (), decreases from 2 to 1. When the system enters the two-phase region, it is no longer possible to independently control temperature and pressure.
In the phase diagram to the right, the boundary curve between the liquid and gas regions maps the constraint between temperature and pressure when the single-component system has separated into liquid and gas phases at equilibrium. The only way to increase the pressure on the two phase line is by increasing the temperature. If the temperature is decreased by cooling, some of the gas condenses, decreasing the pressure. Throughout both processes, the temperature and pressure stay in the relationship shown by this boundary curve unless one phase is entirely consumed by evaporation or condensation, or unless the critical point is reached. As long as there are two phases, there is only one degree of freedom, which corresponds to the position along the phase boundary curve.
The critical point is the black dot at the end of the liquid–gas boundary. As this point is approached, the liquid and gas phases become progressively more similar until, at the critical point, there is no longer a separation into two phases. Above the critical point and away from the phase boundary curve, and the temperature and pressure can be controlled independently. Hence there is only one phase, and it has the physical properties of a dense gas, but is also referred to as a supercritical fluid.
Of the other two-boundary curves, one is the solid–liquid boundary or melting point curve which indicates the conditions for equilibrium between these two phases, and the other at lower temperature and pressure is the solid–gas boundary.
Even for a pure substance, it is possible that three phases, such as solid, liquid and vapour, can exist together in equilibrium (). If there is only one component, there are no degrees of freedom () when there are three phases. Therefore, in a single-component system, this three-phase mixture can only exist at a single temperature and pressure, which is known as a triple point. Here there are two equations , which are sufficient to determine the two variables T and p. In the diagram for CO2 the triple point is the point at which the solid, liquid and gas phases come together, at 5.2 bar and 217 K. It is also possible for other sets of phases to form a triple point, for example in the water system there is a triple point where ice I, ice III and liquid can coexist.
If four phases of a pure substance were in equilibrium (), the phase rule would give , which is meaningless, since there cannot be −1 independent variables. This explains the fact that four phases of a pure substance (such as ice I, ice III, liquid water and water vapour) are not found in equilibrium at any temperature and pressure. In terms of chemical potentials there are now three equations, which cannot in general be satisfied by any values of the two variables T and p, although in principle they might be solved in a special case where one equation is mathematically dependent on the other two. In practice, however, the coexistence of more phases than allowed by the phase rule normally means that the phases are not all in true equilibrium.
Two-component systems
For binary mixtures of two chemically independent components, so that . In addition to temperature and pressure, the other degree of freedom is the composition of each phase, often expressed as mole fraction or mass fraction of one component.
As an example, consider the system of two completely miscible liquids such as toluene and benzene, in equilibrium with their vapours. This system may be described by a boiling-point diagram which shows the composition (mole fraction) of the two phases in equilibrium as functions of temperature (at a fixed pressure).
Four thermodynamic variables which may describe the system include temperature (T), pressure (p), mole fraction of component 1 (toluene) in the liquid phase (x1L), and mole fraction of component 1 in the vapour phase (x1V). However, since two phases are present () in equilibrium, only two of these variables can be independent (). This is because the four variables are constrained by two relations: the equality of the chemical potentials of liquid toluene and toluene vapour, and the corresponding equality for benzene.
For given T and p, there will be two phases at equilibrium when the overall composition of the system (system point) lies in between the two curves. A horizontal line (isotherm or tie line) can be drawn through any such system point, and intersects the curve for each phase at its equilibrium composition. The quantity of each phase is given by the lever rule (expressed in the variable corresponding to the x-axis, here mole fraction).
For the analysis of fractional distillation, the two independent variables are instead considered to be liquid-phase composition (x1L) and pressure. In that case the phase rule implies that the equilibrium temperature (boiling point) and vapour-phase composition are determined.
Liquid–vapour phase diagrams for other systems may have azeotropes (maxima or minima) in the composition curves, but the application of the phase rule is unchanged. The only difference is that the compositions of the two phases are equal exactly at the azeotropic composition.
Aqueous solution of 4 kinds of salts
Consider an aqueous solution containing sodium chloride (NaCl), potassium chloride (KCl), sodium bromide (NaBr), and potassium bromide (KBr), in equilibrium with their respective solid phases. Each salt, in solid form, is a different phase, because each possesses a distinct crystal structure and composition. The aqueous solution itself is another phase, because it forms a homogeneous liquid phase separate from the solid salts, with its own distinct composition and physical properties. Thus we have P = 5 phases.
There are 6 elements present (H, O, Na, K, Cl, Br), but we have 2 constraints:
The stoichiometry of water: n(H) = 2n(O).
Charge balance in the solution: n(Na) + n(K) = n(Cl) + n(Br).
giving C = 6 - 2 = 4 components. The Gibbs phase rule states that F = 1. So, for example, if we plot the P-T phase diagram of the system, there is only one line at which all phases coexist. Any deviation from the line would either cause one of the salts to completely dissolve or one of the ions to completely precipitate from the solution.
Phase rule at constant pressure
For applications in materials science dealing with phase changes between different solid structures, pressure is often imagined to be constant (for example at 1 atmosphere), and is ignored as a degree of freedom, so the formula becomes:
This is sometimes incorrectly called the "condensed phase rule", but it is not applicable to condensed systems subject to high pressures (for example, in geology), since the effects of these pressures are important.
Phase rule in colloidal mixtures
In colloidal mixtures quintuple and sixtuple points have been described in violation of Gibbs phase rule but it is argued that in these systems the rule can be generalized to where accounts for additional parameters of interaction among the components like the diameter of one type of particle in relation to the diameter of the other particles in the solution.
| Physical sciences | Phase transitions | Physics |
304011 | https://en.wikipedia.org/wiki/Onyx | Onyx | Onyx is the parallel-banded variety of chalcedony, a silicate mineral. Agate and onyx are both varieties of layered chalcedony that differ only in the form of the bands. Onyx has parallel bands, while agate has curved bands. The colors of its bands range from black to almost every color. Specimens of onyx commonly contain bands of black or white or both. Onyx, as a descriptive term, has also been applied to parallel-banded varieties of alabaster, marble, calcite, obsidian, and opal, and misleadingly to materials with contorted banding, such as "cave onyx" and "Mexican onyx".
Etymology
Onyx comes through Latin (of the same spelling), from the Ancient Greek (), meaning or . Onyx with pink and white bands can sometimes resemble a fingernail. The English word "nail" is cognate with the Greek word.
Varieties
Onyx is formed of chalcedony bands in alternating colors. It is cryptocrystalline, consisting of fine intergrowths of the silica minerals quartz and moganite. Its bands are parallel, unlike the more chaotic banding that often occurs in agates.
Sardonyx is a variant in which the colored bands are sard (shades of red) rather than black. Black onyx is perhaps the most famous variety, but it is not as common as onyx with colored bands. Artificial treatments have been used since ancient times to produce the black color in "black onyx" and the reds and yellows in sardonyx. Most "black onyx" on the market is artificially colored.
Imitations and treatments
The name has also commonly been used to label other banded materials, such as banded calcite found in Mexico, India, and other places, and often carved, polished, and sold. This material is much softer than true onyx and more readily available. The majority of carved items sold as "onyx" today are this carbonate material.
Artificial onyx types have also been produced from common chalcedony and plain agates. The first-century naturalist Pliny the Elder described these techniques used in Roman times. Treatments for producing black and other colors include soaking or boiling chalcedony in sugar solutions, then treating with sulfuric or hydrochloric acid to carbonize sugars which had been absorbed into the top layers of the stone. These techniques are still used, as well as other dyeing treatments, and most so-called "black onyx" sold is artificially treated. In addition to dye treatments, heating and treatment with nitric acid have been used to lighten or eliminate undesirable colors.
Geographic occurrence
Onyx can be found in various regions of the world, including Greece, Yemen, Uruguay, Argentina, Australia, Brazil, Canada, China, Czech Republic, Germany, Pakistan, India, Indonesia, Madagascar, Latin America, the UK, and various states in the US.
Historical use
It has a long history of use for hardstone carving and jewelry, where it is usually cut as a cabochon or into beads. It has also been used for intaglio and hardstone cameo engraved gems, where the bands make the image contrast with the ground. Some onyx is natural but much of the material in commerce is produced by the staining of agate.
Onyx was used in Egypt as early as the Second Dynasty to make bowls and other pottery items. Use of sardonyx appears in the art of Minoan Crete, notably from the archaeological recoveries at Knossos.
Brazilian green onyx was often used as plinths for art deco sculptures created in the 1920s and 1930s. The German sculptor Ferdinand Preiss used Brazilian green onyx for the base on the majority of his chryselephantine sculptures. Green onyx was also used for trays and pin dishes – produced mainly in Austria – often with small bronze animals or figures attached.
Onyx is mentioned in the Bible many times. Sardonyx (onyx in which white layers alternate with sard - a brownish color) is mentioned in the Bible as well.
Onyx was known to the Ancient Greeks and Romans. The first-century naturalist Pliny the Elder described both types of onyx and various artificial treatment techniques in his Naturalis Historia.
Slabs of onyx (from the Atlas Mountains) were famously used by Mies van der Rohe in Villa Tugendhat at Brno (completed 1930) to create a shimmering semi-translucent interior wall.
The Hôtel de la Païva in Paris is noted for its yellow onyx décor, and the new Mariinsky Theatre Second Stage in St.Petersburg uses yellow onyx in the lobby.
Superstitions
The ancient Romans entered battle carrying amulets of sardonyx engraved with Mars, the god of war. This was believed to bestow courage in battle. In Renaissance Europe, wearing sardonyx was believed to bestow eloquence. A traditional Persian belief is that it helped with epilepsy. Sardonyx was traditionally used by English midwives to ease childbirth by laying it between the breasts of the mother.
| Physical sciences | Silicate minerals | Earth science |
304091 | https://en.wikipedia.org/wiki/Annatto | Annatto | Annatto ( or ) is an orange-red condiment and food coloring derived from the seeds of the achiote tree (Bixa orellana), native to tropical parts of the Americas. It is often used to impart a yellow to red-orange color to foods, but sometimes also for its flavor and aroma. Its scent is described as "slightly peppery with a hint of nutmeg" and flavor as "slightly nutty, sweet and peppery".
The color of annatto comes from various carotenoid pigments, mainly bixin and norbixin, found in the reddish waxy coating of the seeds. The condiment is typically prepared by grinding the seeds to a powder or paste. Similar effects can be obtained by extracting some of the color and flavor principles from the seeds with hot water, oil, or lard, which are then added to the food.
Annatto and its extracts are now widely used in an artisanal or industrial scale as a coloring agent in many processed food products, such as cheeses, dairy spreads, butter and margarine, custards, cakes and other baked goods, potatoes, snack foods, breakfast cereals, smoked fish, sausages, and more. In these uses, annatto is a natural alternative to synthetic food coloring compounds, but it has been linked to rare cases of food-related allergies. Annatto is of particular commercial value in the United States because the Food and Drug Administration considers colorants derived from it to be "exempt of certification".
History
The annatto tree B. orellana is believed to originate in tropical regions from Mexico to Brazil. It was probably not initially used as a food additive, but for other purposes, such as ritual and decorative body painting (still an important tradition in many Brazilian native tribes, such as the Wari'); sunscreen; insect repellent; and for medical purposes. It was used for Mexican manuscript painting in the 16th century. Men of the Tsàchila tribe in Ecuador are highly recognizable thanks to their traditional bright orange hair, which is achieved by using crushed seeds of annatto. It is believed they have been doing so for centuries.
Annatto has been traditionally used as both a coloring and flavoring agent in various cuisines from Latin America, the Caribbean, the Philippines, and other countries where it was taken home by Spanish and Portuguese colonizers in the 16th century. It has various local names according to region. Its use has spread in historic times to other parts of the world, and it was incorporated in local culinary traditions of many countries outside the Americas.
Culinary uses
Traditional cuisine
Ground annatto seeds, often mixed with other seeds or spices, are used in the form of paste or powder for culinary use, especially in Latin American, Jamaican, Belizean, Chamorro, Vietnamese, and Filipino cuisines. In Mexican and Belizean cuisines, it is used to make the spice recado rojo. In Venezuela, annatto is used in the preparation of hallacas, huevos pericos, and other traditional dishes. In Puerto Rico, it is often simmered in oil or ground with seasonings and herbs to make sazón or used to make pasteles, arroz con gandules, and several other dishes, where it is one of the main ingredients. Annatto paste is an important ingredient of cochinita pibil, the slow-roasted pork dish popular in Mexico. It is also a key ingredient in the drink tascalate from Chiapas, Mexico. In the Philippines, it is used for the sauce of pancit palabok. In Guam, it is used to make a staple rice dish flavored with annatto, onion, garlic, butter, and other spices.
In earlier times, the seeds of the plant were also used as a staple food in South and Central America. Today, an Inca recipe for beer made from cocoa, chilli, honey and annatto is available again. The recipe dates back to 1200 BC and is now marketed in North America under the trade name Dogfish Head Theobroma. In the prehistory and early history of South and Central America, the pure, soaked annatto seed was fermented with 5% honey for 10 days and sometimes used as a sole foodstuff. The continent's universities have now recognised the high nutrient content of the seed and are researching ways to reintroduce annatto as a foodstuff. However, practical solutions that would guarantee a basic food supply from the wild harvest have not yet been taken up again by the population.
Industrial food coloring
Annatto is commonly used to impart a yellow or orange color to many industrialized and semi-industrialized foods, including cheese, ice cream, bakery products, desserts, fruit fillings, yogurt, butter, oils, margarines, processed cheese, and fat-based products. In the United States, annatto extract is listed as a color additive "exempt from certification" and is informally considered to be a natural coloring. Foods colored with annatto may declare the coloring in the statement of ingredients as "colored with annatto" or "annatto color". In the European Union, it is identified by the E number E160b.
Cheese
In cheese, the yellow and orange hues naturally vary throughout the year as the cow's feed changes: in the summer, with fresh grass and its natural carotene content, the milk produced would have a natural orange tint, as would the cheese made from it, while at other times of the year, the tint would be greatly reduced. As the pigment is carried in the cream, skimming the milk, which some farmers did to make butter or to sell it separately, the lesser-quality cheese from such milk would be white.
To fool the consumer, the cheesemakers introduced colorants to imitate the more intense colors of the finer summer cheese. Initially these colors came from saffron, marigold, and carrot juice, but later annatto began being used.
In the 17th century, the Dutch, who had established colonies in Guyana, traded in food, particularly an orange-red natural colorant, annatto, with the indigenous communities. Zeeland traders under the authority of the West India Company bought annatto from the inhabitants of the coastal regions of Guyana and Suriname and sold it in the Netherlands as verw ('paint'). One contemporaneous description comes from Adriaen van Berkel, in a book published in 1695, though he does not mention whether it was used in cheese.
The earliest known documentation of annatto's use in cheese is in a 1743 Dutch volume (Household Dictionary), according to American scientist Paul Kindstedt of the University of Vermont. Other historical documents from the period confirm that annatto (then called "orleaan" or "orleans") was being used to color cheese by the mid-18th century.
England is another country that has used annatto to color its cheeses; colorants have been added to Gloucester cheese as early as the 16th century to allow inferior cheese to masquerade as the best Double Gloucester, with annatto later being used for that purpose. This usage was subsequently adopted in other parts of the UK, for cheeses such as Cheshire and Red Leicester, as well as colored Cheddar made in Scotland. Many cheddars are produced in both white and red (orange) varieties, the only difference between the two being the presence of annatto as a coloring. That practice has extended to many modern processed cheese products, such as American cheese and Velveeta. Cheeses from other countries also use annatto, including Mimolette from France and Leyden from the Netherlands.
Cheeses that use annatto in at least some preparations include:
Chemical composition
The yellow to orange color is produced by the chemical compounds bixin and norbixin, which are classified as carotenoids. The fat-soluble color in the crude extract is called bixin, which can then be saponified into water-soluble norbixin. This dual solubility property of annatto is rare for carotenoids. The seeds contain 4.5–5.5% pigment, which consists of 70–80% bixin. Unlike beta-carotene, another well-known carotenoid, annatto-based pigments are not vitamin A precursors. The more norbixin in an annatto preparation, the more yellow it is; a higher level of bixin gives it a more orange hue.
Safety
Annatto condiments and colorants are safe for most people when used in food amounts, but they may cause allergic reactions in those who are sensitive. In one 1978 study of 61 patients with chronic hives or angioedema, 56 patients were orally provoked by annatto extract during an elimination diet. A challenge was performed with a dose equivalent to the amount used in of butter. Twenty-six percent of the patients reacted to this color four hours after intake, worse than synthetic dyes, such as amaranth (9%), tartrazine (11%), sunset yellow FCF (17%), allura red AC (16%), ponceau 4R (15%), erythrosine (12%) and brilliant blue FCF (14%).
Annatto is not among the "Big Eight" substances causing hypersensitivity reactions which are responsible for more than 90% of food related allergic reactions. The US FDA and experts at the Food Allergy Research and Resource Program (FARRP) of the University of Nebraska do not include annatto in the list of major food allergens.
| Biology and health sciences | Herbs and spices | Plants |
304471 | https://en.wikipedia.org/wiki/Case%20study | Case study | A case study is an in-depth, detailed examination of a particular case (or cases) within a real-world context. For example, case studies in medicine may focus on an individual patient or ailment; case studies in business might cover a particular firm's strategy or a broader market; similarly, case studies in politics can range from a narrow happening over time like the operations of a specific political campaign, to an enormous undertaking like world war, or more often the policy analysis of real-world problems affecting multiple stakeholders.
Generally, a case study can highlight nearly any individual, group, organization, event, belief system, or action. A case study does not necessarily have to be one observation (N=1), but may include many observations (one or multiple individuals and entities across multiple time periods, all within the same case study). Research projects involving numerous cases are frequently called cross-case research, whereas a study of a single case is called within-case research.
Case study research has been extensively practiced in both the social and natural sciences.
Definition
There are multiple definitions of case studies, which may emphasize the number of observations (a small N), the method (qualitative), the thickness of the research (a comprehensive examination of a phenomenon and its context), and the naturalism (a "real-life context" is being examined) involved in the research. There is general agreement among scholars that a case study does not necessarily have to entail one observation (N=1), but can include many observations within a single case or across numerous cases. For example, a case study of the French Revolution would at the bare minimum be an observation of two observations: France before and after a revolution. John Gerring writes that the N=1 research design is so rare in practice that it amounts to a "myth".
The term cross-case research is frequently used for studies of multiple cases, whereas within-case research is frequently used for a single case study.
John Gerring defines the case study approach as an "intensive study of a single unit or a small number of units (the cases), for the purpose of understanding a larger class of similar units (a population of cases)". According to Gerring, case studies lend themselves to an idiographic style of analysis, whereas quantitative work lends itself to a nomothetic style of analysis. He adds that "the defining feature of qualitative work is its use of noncomparable observations—observations that pertain to different aspects of a causal or descriptive question", whereas quantitative observations are comparable.
According to John Gerring, the key characteristic that distinguishes case studies from all other methods is the "reliance on evidence drawn from a single case and its attempts, at the same time, to illuminate features of a broader set of cases". Scholars use case studies to shed light on a "class" of phenomena.
Research designs
As with other social science methods, no single research design dominates case study research. Case studies can use at least four types of designs. First, there may be a "no theory first" type of case study design, which is closely connected to Kathleen M. Eisenhardt's methodological work. A second type of research design highlights the distinction between single- and multiple-case studies, following Robert K. Yin's guidelines and extensive examples. A third design deals with a "social construction of reality", represented by the work of Robert E. Stake. Finally, the design rationale for a case study may be to identify "anomalies". A representative scholar of this design is Michael Burawoy. Each of these four designs may lead to different applications, and understanding their sometimes unique ontological and epistemological assumptions becomes important. However, although the designs can have substantial methodological differences, the designs also can be used in explicitly acknowledged combinations with each other.
While case studies can be intended to provide bounded explanations of single cases or phenomena, they are often intended to raise theoretical insights about the features of a broader population.
Case selection and structure
Case selection in case study research is generally intended to find cases that are representative samples and which have variations on the dimensions of theoretical interest. Using that is solely representative, such as an average or typical case is often not the richest in information. In clarifying lines of history and causation it is more useful to select subjects that offer an interesting, unusual, or particularly revealing set of circumstances. A case selection that is based on representativeness will seldom be able to produce these kinds of insights.
While a random selection of cases is a valid case selection strategy in large-N research, there is a consensus among scholars that it risks generating serious biases in small-N research. Random selection of cases may produce unrepresentative cases, as well as uninformative cases. Cases should generally be chosen that have a high expected information gain. For example, outlier cases (those which are extreme, deviant or atypical) can reveal more information than the potentially representative case. A case may also be chosen because of the inherent interest of the case or the circumstances surrounding it. Alternatively, it may be chosen because of researchers' in-depth local knowledge; where researchers have this local knowledge they are in a position to "soak and poke" as Richard Fenno put it, and thereby to offer reasoned lines of explanation based on this rich knowledge of setting and circumstances.
Beyond decisions about case selection and the subject and object of the study, decisions need to be made about the purpose, approach, and process of the case study. Gary Thomas thus proposes a typology for the case study wherein purposes are first identified (evaluative or exploratory), then approaches are delineated (theory-testing, theory-building, or illustrative), then processes are decided upon, with a principal choice being between whether the study is to be single or multiple, and choices also about whether the study is to be retrospective, snapshot or diachronic, and whether it is nested, parallel or sequential.
In a 2015 article, John Gerring and Jason Seawright list seven case selection strategies:
Typical cases are cases that exemplify a stable cross-case relationship. These cases are representative of the larger population of cases, and the purpose of the study is to look within the case rather than compare it with other cases.
Diverse cases are cases that have variations on the relevant X and Y variables. Due to the range of variation on the relevant variables, these cases are representative of the full population of cases.
Extreme cases are cases that have an extreme value on the X or Y variable relative to other cases.
Deviant cases are cases that defy existing theories and common sense. They not only have extreme values on X or Y (like extreme cases) but defy existing knowledge about causal relations.
Influential cases are cases that are central to a model or theory (for example, Nazi Germany in theories of fascism and the far-right).
Most similar cases are cases that are similar on all the independent variables, except the one of interest to the researcher.
Most different cases are cases that are different on all the independent variables, except the one of interest to the researcher.
For theoretical discovery, Jason Seawright recommends using deviant cases or extreme cases that have an extreme value on the X variable.
Arend Lijphart, and Harry Eckstein identified five types of case study research designs (depending on the research objectives), Alexander George and Andrew Bennett added a sixth category:
Atheoretical (or configurative idiographic) case studies aim to describe a case very well, but not to contribute to a theory.
Interpretative (or disciplined configurative) case studies aim to use established theories to explain a specific case.
Hypothesis-generating (or heuristic) case studies aim to inductively identify new variables, hypotheses, causal mechanisms, and causal paths.
Theory testing case studies aim to assess the validity and scope conditions of existing theories.
Plausibility probes, aim to assess the plausibility of new hypotheses and theories.
Building block studies of types or subtypes, aim to identify common patterns across cases.
Aaron Rapport reformulated "least-likely" and "most-likely" case selection strategies into the "countervailing conditions" case selection strategy. The countervailing conditions case selection strategy has three components:
The chosen cases fall within the scope conditions of both the primary theory being tested and the competing alternative hypotheses.
For the theories being tested, the analyst must derive clearly stated expected outcomes.
In determining how difficult a test is, the analyst should identify the strength of countervailing conditions in the chosen cases.
In terms of case selection, Gary King, Robert Keohane, and Sidney Verba warn against "selecting on the dependent variable". They argue for example that researchers cannot make valid causal inferences about war outbreaks by only looking at instances where war did happen (the researcher should also look at cases where war did not happen). Scholars of qualitative methods have disputed this claim, however. They argue that selecting the dependent variable can be useful depending on the purposes of the research. Barbara Geddes shares their concerns with selecting the dependent variable (she argues that it cannot be used for theory testing purposes), but she argues that selecting on the dependent variable can be useful for theory creation and theory modification.
King, Keohane, and Verba argue that there is no methodological problem in selecting the explanatory variable, however. They do warn about multicollinearity (choosing two or more explanatory variables that perfectly correlate with each other).
Uses
Case studies have commonly been seen as a fruitful way to come up with hypotheses and generate theories. Case studies are useful for understanding outliers or deviant cases. Classic examples of case studies that generated theories includes Darwin's theory of evolution (derived from his travels to the Easter Island), and Douglass North's theories of economic development (derived from case studies of early developing states, such as England).
Case studies are also useful for formulating concepts, which are an important aspect of theory construction. The concepts used in qualitative research will tend to have higher conceptual validity than concepts used in quantitative research (due to conceptual stretching: the unintentional comparison of dissimilar cases). Case studies add descriptive richness, and can have greater internal validity than quantitative studies. Case studies are suited to explain outcomes in individual cases, which is something that quantitative methods are less equipped to do.
Case studies have been characterized as useful to assess the plausibility of arguments that explain empirical regularities. Case studies are also useful for understanding outliers or deviant cases.
Through fine-gained knowledge and description, case studies can fully specify the causal mechanisms in a way that may be harder in a large-N study. In terms of identifying "causal mechanisms", some scholars distinguish between "weak" and "strong chains". Strong chains actively connect elements of the causal chain to produce an outcome whereas weak chains are just intervening variables.
Case studies of cases that defy existing theoretical expectations may contribute knowledge by delineating why the cases violate theoretical predictions and specifying the scope conditions of the theory. Case studies are useful in situations of causal complexity where there may be equifinality, complex interaction effects and path dependency. They may also be more appropriate for empirical verifications of strategic interactions in rationalist scholarship than quantitative methods. Case studies can identify necessary and insufficient conditions, as well as complex combinations of necessary and sufficient conditions. They argue that case studies may also be useful in identifying the scope conditions of a theory: whether variables are sufficient or necessary to bring about an outcome.
Qualitative research may be necessary to determine whether a treatment is as-if random or not. As a consequence, good quantitative observational research often entails a qualitative component.
Limitations
Designing Social Inquiry (also called "KKV"), an influential 1994 book written by Gary King, Robert Keohane, and Sidney Verba, primarily applies lessons from regression-oriented analysis to qualitative research, arguing that the same logics of causal inference can be used in both types of research. The authors' recommendation is to increase the number of observations (a recommendation that Barbara Geddes also makes in Paradigms and Sand Castles), because few observations make it harder to estimate multiple causal effects, as well as increase the risk that there is measurement error, and that an event in a single case was caused by random error or unobservable factors. KKV sees process-tracing and qualitative research as being "unable to yield strong causal inference" due to the fact that qualitative scholars would struggle with determining which of many intervening variables truly links the independent variable with a dependent variable. The primary problem is that qualitative research lacks a sufficient number of observations to properly estimate the effects of an independent variable. They write that the number of observations could be increased through various means, but that would simultaneously lead to another problem: that the number of variables would increase and thus reduce degrees of freedom. Christopher H. Achen and Duncan Snidal similarly argue that case studies are not useful for theory construction and theory testing.
The purported "degrees of freedom" problem that KKV identify is widely considered flawed; while quantitative scholars try to aggregate variables to reduce the number of variables and thus increase the degrees of freedom, qualitative scholars intentionally want their variables to have many different attributes and complexity. For example, James Mahoney writes, "the Bayesian nature of process of tracing explains why it is inappropriate to view qualitative research as suffering from a small-N problem and certain standard causal identification problems." By using Bayesian probability, it may be possible to makes strong causal inferences from a small sliver of data.
KKV also identify inductive reasoning in qualitative research as a problem, arguing that scholars should not revise hypotheses during or after data has been collected because it allows for ad hoc theoretical adjustments to fit the collected data. However, scholars have pushed back on this claim, noting that inductive reasoning is a legitimate practice (both in qualitative and quantitative research).
A commonly described limit of case studies is that they do not lend themselves to generalizability. Due to the small number of cases, it may be harder to ensure that the chosen cases are representative of the larger population.
As small-N research should not rely on random sampling, scholars must be careful in avoiding selection bias when picking suitable cases. A common criticism of qualitative scholarship is that cases are chosen because they are consistent with the scholar's preconceived notions, resulting in biased research. Alexander George and Andrew Bennett also note that a common problem in case study research is that of reconciling conflicting interpretations of the same data. Another limit of case study research is that it can be hard to estimate the magnitude of causal effects.
Teaching case studies
Teachers may prepare a case study that will then be used in classrooms in the form of a "teaching" case study (also see case method and casebook method). For instance, as early as 1870 at Harvard Law School, Christopher Langdell departed from the traditional lecture-and-notes approach to teaching contract law and began using cases pled before courts as the basis for class discussions. By 1920, this practice had become the dominant pedagogical approach used by law schools in the United States.
Outside of law, teaching case studies have become popular in many different fields and professions, ranging from business education to science education. The Harvard Business School has been among the most prominent developers and users of teaching case studies. Teachers develop case studies with particular learning objectives in mind. Additional relevant documentation, such as financial statements, time-lines, short biographies, and multimedia supplements (such as video-recordings of interviews) often accompany the case studies. Similarly, teaching case studies have become increasingly popular in science education, covering different biological and physical sciences. The National Center for Case Studies in Teaching Science has made a growing body of teaching case studies available for classroom use, for university as well as secondary school coursework.
| Physical sciences | Research methods | Basics and measurement |
304588 | https://en.wikipedia.org/wiki/Brain%20injury | Brain injury | Brain injury (BI) is the destruction or degeneration of brain cells. Brain injuries occur due to a wide range of internal and external factors. In general, brain damage refers to significant, undiscriminating trauma-induced damage.
A common category with the greatest number of injuries is traumatic brain injury (TBI) following physical trauma or head injury from an outside source, and the term acquired brain injury (ABI) is used in appropriate circles to differentiate brain injuries occurring after birth from injury, from a genetic disorder (GBI), or from a congenital disorder (CBI). Primary and secondary brain injuries identify the processes involved, while focal and diffuse brain injury describe the severity and localization.
Impaired function of affected areas can be compensated through neuroplasticity by forming new neural connections.
Signs and symptoms
Symptoms of brain injuries vary based on the severity of the injury or how much of the brain is affected. The four categories used for classifying the severity of brain injuries are mild, moderate, or severe.
Severity of injuries
Mild brain injuries
Symptoms of a mild brain injury include headaches, confusions, tinnitus, fatigue, changes in sleep patterns, mood or behavior. Other symptoms include trouble with memory, concentration, attention or thinking. Mental fatigue is a common debilitating experience and may not be linked by the patient to the original (minor) incident.
Moderate/severe brain injuries
Cognitive symptoms include confusion, aggressiveness, abnormal behavior, slurred speech, and coma or other disorders of consciousness. Physical symptoms include headaches that worsen or do not go away, vomiting or nausea, convulsions, brain pulsation, abnormal dilation of the eyes, inability to awaken from sleep, weakness in extremities, and loss of coordination.
Symptoms in children
Symptoms observed in children include changes in eating habits, persistent irritability or sadness, changes in attention, or disrupted sleeping habits.
Location of brain damage predicts symptoms
Symptoms of brain injuries can also be influenced by the location of the injury and as a result impairments are specific to the part of the brain affected. Lesion size is correlated with severity, recovery, and comprehension. Brain injuries often create impairment or disability that can vary greatly in severity.
In cases of severe brain injuries, the likelihood of areas with permanent disability is great, including neurocognitive deficits, delusions (often, to be specific, monothematic delusions), speech or movement problems, and intellectual disability. There may also be personality changes. The most severe cases result in coma or even persistent vegetative state. Even a mild incident can have long-term effects or cause symptoms to appear years later.
Studies show there is a correlation between brain lesion and language, speech, and category-specific disorders. Wernicke's aphasia is associated with anomia, unknowingly making up words (neologisms), and problems with comprehension. The symptoms of Wernicke's aphasia are caused by damage to the posterior section of the superior temporal gyrus.
Damage to the Broca's area typically produces symptoms like omitting functional words (agrammatism), sound production changes, dyslexia, dysgraphia, and problems with comprehension and production. Broca's aphasia is indicative of damage to the posterior inferior frontal gyrus of the brain.
An impairment following damage to a region of the brain does not necessarily imply that the damaged area is wholly responsible for the cognitive process which is impaired, however. For example, in pure alexia, the ability to read is destroyed by a lesion damaging both the left visual field and the connection between the right visual field and the language areas (Broca's area and Wernicke's area). However, this does not mean one with pure alexia is incapable of comprehending speech—merely that there is no connection between their working visual cortex and language areas—as is demonstrated by the fact that people with pure alexia can still write, speak, and even transcribe letters without understanding their meaning.
Lesions to the fusiform gyrus often result in prosopagnosia, the inability to distinguish faces and other complex objects from each other. Lesions in the amygdala would eliminate the enhanced activation seen in occipital and fusiform visual areas in response to fear with the area intact. Amygdala lesions change the functional pattern of activation to emotional stimuli in regions that are distant from the amygdala.
Other lesions to the visual cortex have different effects depending on the location of the damage. Lesions to V1, for example, can cause blindsight in different areas of the brain depending on the size of the lesion and location relative to the calcarine fissure. Lesions to V4 can cause color-blindness, and bilateral lesions to MT/V5 can cause the loss of the ability to perceive motion. Lesions to the parietal lobes may result in agnosia, an inability to recognize complex objects, smells, or shapes, or amorphosynthesis, a loss of perception on the opposite side of the body.
Non-localizing features
Brain injuries have far-reaching and varied consequences due to the nature of the brain as the main source of bodily control. Brain-injured people commonly experience issues with memory. This can be issues with either long or short-term memories depending on the location and severity of the injury. Sometimes memory can be improved through rehabilitation, although it can be permanent. Behavioral and personality changes are also commonly observed due to changes of the brain structure in areas controlling hormones or major emotions.
Headaches and pain can occur as a result of a brain injury, either directly from the damage or due to neurological conditions stemming from the injury. Due to the changes in the brain as well as the issues associated with the change in physical and mental capacity, depression and low self-esteem are common side effects that can be treated with psychological help. Antidepressants must be used with caution in brain injury people due to the potential for undesired effects because of the already altered brain chemistry.
Long term psychological and physiological effects
There are multiple responses of the body to brain injury, occurring at different times after the initial occurrence of damage, as the functions of the neurons, nerve tracts, or sections of the brain can be affected by damage. The immediate response can take many forms. Initially, there may be symptoms such as swelling, pain, bruising, or loss of consciousness. Post-traumatic amnesia is also common with brain damage, as is temporary aphasia, or impairment of language.
As time progresses, and the severity of injury becomes clear, there are further responses that may become apparent. Due to loss of blood flow or damaged tissue, sustained during the injury, amnesia and aphasia may become permanent, and apraxia has been documented in patients. Amnesia is a condition in which a person is unable to remember things. Aphasia is the loss or impairment of word comprehension or use. Apraxia is a motor disorder caused by damage to the brain, and may be more common in those who have been left brain damaged, with loss of mechanical knowledge critical. Headaches, occasional dizziness, and fatigue—all temporary symptoms of brain trauma—may become permanent, or may not disappear for a long time.
There are documented cases of lasting psychological effects as well, such as emotional changes often caused by damage to the various parts of the brain that control human emotions and behavior. Individuals who have experienced emotional changes related to brain damage may have emotions that come very quickly and are very intense, but have very little lasting effect. Emotional changes may not be triggered by a specific event, and can be a cause of stress to the injured party and their family and friends. Often, counseling is suggested for those who experience this effect after their injury, and may be available as an individual or group session.
The long term psychological and physiological effects will vary by person and injury. For example, perinatal brain damage has been implicated in cases of neurodevelopmental impairments and psychiatric illnesses. If any concerning symptoms, signs, or changes to behaviors are occurring, a healthcare provider should be consulted.
Causes
Brain injuries can result from a number of conditions, including:
Trauma; multiple traumatic injuries can lead to chronic traumatic encephalopathy. A coup-contrecoup injury occurs when the force impacting the head is not only strong enough to cause a contusion at the site of impact, but also able to move the brain and cause it to displace rapidly into the opposite side of the skull, causing an additional contusion.
open head injury
closed head injury
penetrating: when a sharp object enters the brain, causing a large damage area. Penetrating injuries caused by bullets have a 91 percent mortality rate.
Deceleration injuries
Poisoning; for example, from heavy metals including mercury and compounds of lead
Genetic disorder
Hypoxia, including birth hypoxia
Tumors
Infections
Stroke leading to infarct, which may follow thrombosis, embolisms, angiomas, aneurysms, and cerebral arteriosclerosis.
Neurological illness or disorders such as cerebral palsy, Parkinson's disease, etc.
Surgery
Substance use disorder
Neurotoxins- pollution exposure or biological exposure (Annonaceae, rotenone, Aspergillus spores, West Nile fever, Viral meningitis)
Suicide attempt such as hanging, falling off from height, and even on rare occasion getting shot by a firearm, etc.
Acute total or REM sleep deprivation lasting longer than a day
Chemotherapy
Chemotherapy can cause brain damage to the neural stem cells and oligodendrocyte cells that produce myelin. Radiation and chemotherapy can lead to brain tissue damage by disrupting or stopping blood flow to the affected areas of the brain. This damage can cause long term effects such as but not limited to; memory loss, confusion, and loss of cognitive function. The brain damage caused by radiation depends on where the brain tumor is located, the amount of radiation used, and the duration of the treatment. Radiosurgery can also lead to tissue damage that results in about 1 in 20 patients requiring a second operation to remove the damaged tissue.
Wernicke–Korsakoff syndrome
Wernicke–Korsakoff syndrome can cause brain damage and results from a Vitamin B deficiency (specifically vitamin B1, thiamine). This syndrome presents with two conditions, Wernicke's encephalopathy and Korsakoff psychosis. Typically Wernicke's encephalopathy precedes symptoms of Korsakoff psychosis. Wernicke's encephalopathy results from focal accumulation of lactic acid, causing problems with vision, coordination, and balance.
Korsakoff psychosis typically follows after the symptoms of Wernicke's decrease. Wernicke-Korsakoff syndrome is typically caused by conditions causing thiamine deficiency, such as chronic heavy alcohol use or by conditions that affect nutritional absorption, including colon cancer, eating disorders and gastric bypass.
Iatrogenic
Brain lesions are sometimes intentionally inflicted during neurosurgery, such as the carefully placed brain lesion used to treat epilepsy and other brain disorders. These lesions are induced by excision or by electric shocks (electrolytic lesions) to the exposed brain or commonly by infusion of excitotoxins to specific areas.
Diffuse axonal
Diffuse axonal injury is caused by shearing forces on the brain leading to lesions in the white matter tracts of the brain. These shearing forces are seen in cases where the brain had a sharp rotational acceleration, and is caused by the difference in density between white matter and grey matter.
Body's response to brain injury
Unlike some of the more obvious responses to brain damage, the body also has invisible physical responses which can be difficult to notice. These will generally be identified by a healthcare provider, especially as they are normal physical responses to brain damage. Cytokines are known to be induced in response to brain injury. These have diverse actions that can cause, exacerbate, mediate and/or inhibit cellular injury and repair. TGFβ seems to exert primarily neuroprotective actions, whereas TNFα might contribute to neuronal injury and exert protective effects. IL-1 mediates ischaemic, excitotoxic, and traumatic brain injury, probably through multiple actions on glia, neurons, and the vasculature. Cytokines may be useful in order to discover novel therapeutic strategies. At the current time, they are already in clinical trials.
Diagnosis
Glasgow Coma Scale (GCS) is the most widely used scoring system used to assess the level of severity of a brain injury. This method is based on the objective observations of specific traits to determine the severity of a brain injury. It is based on three traits: eye opening, verbal response, and motor response, gauged as described below. Based on the Glasgow Coma Scale severity is classified as follows, severe brain injuries score 3–8, moderate brain injuries score 9–12 and mild score 13–15.
There are several imaging techniques that can aid in diagnosing and assessing the extent of brain damage, such as computed tomography (CT) scan, magnetic resonance imaging (MRI), diffusion tensor imaging (DTI) magnetic resonance spectroscopy (MRS), positron emission tomography (PET), and single-photon emission tomography (SPECT). CT scans and MRI are the two techniques widely used and are most effective. CT scans can show brain bleeds, fractures of the skull, fluid build up in the brain that will lead to increased cranial pressure.
MRI is able to better to detect smaller injuries, detect damage within the brain, diffuse axonal injury, injuries to the brainstem, posterior fossa, and subtemporal and subfrontal regions. However, patients with pacemakers, metallic implants, or other metal within their bodies are unable to have an MRI done. Typically the other imaging techniques are not used in a clinical setting because of the cost, lack of availability.
Management
Acute
The treatment for emergency traumatic brain injuries focuses on assuring the person has enough oxygen from the brain's blood supply, and on maintaining normal blood pressure to avoid further injuries of the head or neck. The person may need surgery to remove clotted blood or repair skull fractures, for which cutting a hole in the skull may be necessary. Medicines used for traumatic injuries are diuretics, anti-seizure or coma-inducing drugs. Diuretics reduce the fluid in tissues lowering the pressure on the brain. In the first week after a traumatic brain injury, a person may have a risk of seizures, which anti-seizure drugs help prevent. Coma-inducing drugs may be used during surgery to reduce impairments and restore blood flow. Mouse NGF has been licensed in China since 2003 and is used to promote neurological recovery in a range of brain injuries, including intracerebral hemorrhage.
In the case of brain damage from traumatic brain injury, dexamethasone and/or Mannitol may be used.
Chronic
Various professions may be involved in the medical care and rehabilitation of someone with an impairment after a brain injury. Neurologists, neurosurgeons, and physiatrists are physicians specialising in treating brain injury. Neuropsychologists (especially clinical neuropsychologists) are psychologists specialising in understanding the effects of brain injury and may be involved in assessing the severity or creating rehabilitation strategies. Occupational therapists may be involved in running rehabilitation programs to help restore lost function or help re-learn essential skills. Registered nurses, such as those working in hospital intensive care units, are able to maintain the health of the severely brain-injured with constant administration of medication and neurological monitoring, including the use of the Glasgow Coma Scale used by other health professionals to quantify extent of orientation.
Physiotherapists also play a significant role in rehabilitation after a brain injury. In the case of a traumatic brain injury (TBI), physiotherapy treatment during the post-acute phase may include sensory stimulation, serial casting and splinting, fitness and aerobic training, and functional training. Sensory stimulation refers to regaining sensory perception through the use of modalities. There is no evidence to support the efficacy of this intervention. Serial casting and splinting are often used to reduce soft tissue contractures and muscle tone. Evidence based research reveals that serial casting can be used to increase passive range of motion (PROM) and decrease spasticity.
Functional training may also be used to treat patients with TBIs. To date, no studies supports the efficacy of sit to stand training, arm ability training and body weight support systems (BWS). Overall, studies suggest that patients with TBIs who participate in more intense rehabilitation programs will see greater benefits in functional skills. More research is required to better understand the efficacy of the treatments mentioned above.
Other treatments for brain injury can include medication, psychotherapy, neuropsychological rehabilitation, neurotherapy and/or surgery.
Prognosis
Prognosis, or the likely progress of a disorder, depends on the nature, location, and cause of the brain damage (see Traumatic brain injury, Focal and diffuse brain injury, Primary and secondary brain injury).
In general, neuroregeneration can occur in the peripheral nervous system but is much rarer and more difficult to assist in the central nervous system (brain or spinal cord). However, in neural development in humans, areas of the brain can learn to compensate for other damaged areas, and may increase in size and complexity and even change function, just as someone who loses a sense may gain increased acuity in another sense—a process termed neuroplasticity.
There are many misconceptions that revolve around brain injuries and brain damage. One misconception is that if someone has brain damage then they cannot fully recover. Recovery depends a variety of factors; such as severity and location. Testing is done to note severity and location. Not everyone fully heals from brain damage, but it is possible to have a full recovery. Brain injuries are very hard to predict in outcome. Many tests and specialists are needed to determine the likelihood of the prognosis. People with minor brain damage can have debilitating side effects; not just severe brain damage has debilitating effects.
The side-effects of a brain injury depend on location and the body's response to injury. Even a mild concussion can have long term effects that may not resolve. Another misconception is that children heal better from brain damage. Children are at greater risk for injury due to lack of maturity. It makes future development hard to predict. This is because different cortical areas mature at different stages, with some major cell populations and their corresponding cognitive faculties remaining unrefined until early adulthood. In the case of a child with frontal brain injury, for example, the impact of the damage may be undetectable until that child fails to develop normal executive functions in his or her late teens and early twenties.
History
The foundation for understanding human behavior and brain injury can be attributed to the case of Phineas Gage and the famous case studies by Paul Broca. The first case study on Phineas Gage's head injury is one of the most astonishing brain injuries in history. In 1848, Phineas Gage was paving way for a new railroad line when he encountered an accidental explosion of a tamping iron straight through his frontal lobe. Gage observed to be intellectually unaffected but was claimed by some to have exemplified post-injury behavioral deficits.
Ten years later, Paul Broca examined two patients exhibiting impaired speech due to frontal lobe injuries. Broca's first patient lacked productive speech. He saw this as an opportunity to address language localization. It was not until Leborgne, informally known as "tan", died when Broca confirmed the frontal lobe lesion from an autopsy. The second patient had similar speech impairments, supporting his findings on language localization. The results of both cases became a vital verification of the relationship between speech and the left cerebral hemisphere. The affected area is known today as Broca's area and the condition as Broca's aphasia.
A few years later, a German neuroscientist, Carl Wernicke, consulted on a stroke patient. The patient experienced neither speech nor hearing impairments, but had a few brain deficits. These deficits included lacking the ability to comprehend what was spoken to him and the words written down. After his death, Wernicke examined his autopsy that found a lesion located in the left temporal region. This area became known as Wernicke's area. Wernicke later hypothesized the relationship between Wernicke's area and Broca's area, which was proven fact.
| Biology and health sciences | Injury | null |
304604 | https://en.wikipedia.org/wiki/Mechatronics | Mechatronics | Mechatronics engineering, also called mechatronics, is an interdisciplinary branch of engineering that focuses on the integration of mechanical engineering, electrical engineering, electronic engineering and software engineering, and also includes a combination of robotics, computer science, telecommunications, systems, control, automation and product engineering.
As technology advances over time, various subfields of engineering have succeeded in both adapting and multiplying. The intention of mechatronics is to produce a design solution that unifies each of these various subfields. Originally, the field of mechatronics was intended to be nothing more than a combination of mechanics, electrical and electronics, hence the name being a portmanteau of the words "mechanics" and "electronics"; however, as the complexity of technical systems continued to evolve, the definition had been broadened to include more technical areas.
The word mechatronics originated in Japanese-English and was created by Tetsuro Mori, an engineer of Yaskawa Electric Corporation. The word mechatronics was registered as trademark by the company in Japan with the registration number of "46-32714" in 1971. The company later released the right to use the word to the public, and the word began being used globally. Currently the word is translated into many languages and is considered an essential term for advanced automated industry.
Many people treat mechatronics as a modern buzzword synonymous with automation, robotics and electromechanical engineering.
French standard NF E 01-010 gives the following definition: "approach aiming at the synergistic integration of mechanics, electronics, control theory, and computer science within product design and manufacturing, in order to improve and/or optimize its functionality".
History
The word mechatronics was registered as trademark by the company in Japan with the registration number of "46-32714" in 1971. The company later released the right to use the word to the public, and the word began being used globally.
With the advent of information technology in the 1980s, microprocessors were introduced into mechanical systems, improving performance significantly. By the 1990s, advances in computational intelligence were applied to mechatronics in ways that revolutionized the field.
Description
A mechatronics engineer unites the principles of mechanics, electrical, electronics, and computing to generate a simpler, more economical and reliable system.
Engineering cybernetics deals with the question of control engineering of mechatronic systems. It is used to control or regulate such a system (see control theory). Through collaboration, the mechatronic modules perform the production goals and inherit flexible and agile manufacturing properties in the production scheme. Modern production equipment consists of mechatronic modules that are integrated according to a control architecture. The most known architectures involve hierarchy, polyarchy, heterarchy, and hybrid. The methods for achieving a technical effect are described by control algorithms, which might or might not utilize formal methods in their design. Hybrid systems important to mechatronics include production systems, synergy drives,
exploration rovers, automotive subsystems such as anti-lock braking systems and spin-assist, and everyday equipment such as autofocus cameras, video, hard disks, CD players and phones.
Subdisciplines
Mechanical
Mechanical engineering is an important part of mechatronics engineering. It includes the study of mechanical nature of how an object works. Mechanical elements refer to mechanical structure, mechanism, thermo-fluid, and hydraulic aspects of a mechatronics system. The study of thermodynamics, dynamics, fluid mechanics, pneumatics and hydraulics. Mechatronics engineer who works a mechanical engineer can specialize in hydraulics and pneumatics systems, where they can be found working in automobile industries. A mechatronics engineer can also design a vehicle since they have strong mechanical and electronical background. Knowledge of software applications such as computer-aided design and computer aided manufacturing is essential for designing products. Mechatronics covers a part of mechanical syllabus which is widely applied in automobile industry.
Mechatronic systems represent a large part of the functions of an automobile. The control loop formed by sensor—information processing—actuator—mechanical (physical) change is found in many systems. The system size can be very different. The anti-lock braking system (ABS) is a mechatronic system. The brake itself is also one. And the control loop formed by driving control (for example cruise control), engine, vehicle driving speed in the real world and speed measurement is a mechatronic system, too. The great importance of mechatronics for automotive engineering is also evident from the fact that vehicle manufacturers often have development departments with "Mechatronics" in their names.
Electronics and electricals
Electronics and telecommunication engineering specializes in electronics devices and telecom devices of a mechatronics system. A mechatronics engineer specialized in electronics and telecommunications have knowledge of computer hardware devices. The transmission of signal is the main application of this subfield of mechatronics. Where digital and analog systems also forms an important part of mechatronics systems. Telecommunications engineering deals with the transmission of information across a medium.
Electronics engineering is related to computer engineering and electrical engineering. Control engineering has a wide range of electronic applications from the flight and propulsion systems of commercial airplanes to the cruise control present in many modern cars. VLSI designing is important for creating integrated circuits. Mechatronics engineers have deep knowledge of microprocessors, microcontrollers, microchips and semiconductors. The application of mechatronics in electronics manufacturing industry can conduct research and development on consumer electronic devices such as mobile phones, computers, cameras etc. For mechatronics engineers it is necessary to learn operating computer applications such as MATLAB and Simulink for designing and developing electronic products.
Mechatronics engineering is a interdisciplinary course, it includes concepts of both electrical and mechanical systems. A mechatronics engineer engages in designing high power transformers or radio-frequency module transmitters.
Avionics
Avionics is also considered a variant of mechatronics as it combines several fields such as electronics and telecom with aerospace engineering. It is the subdiscipline of mechatronics engineering and aerospace engineering which is engineering branch focusing on electronics systems of aircraft. The word avionics is a blend of aviation and electronics. The electronics system of aircraft includes aircraft communication addressing and reporting system, air navigation, aircraft flight control system, aircraft collision avoidance systems, flight recorder, weather radar and lightning detector. These can be as simple as a searchlight for a police helicopter or as complicated as the tactical system for an airborne early warning platform.
Advanced mechatronics
Another variant is motion control for advanced mechatronics, presently recognized as a key technology in mechatronics. The robustness of motion control will be represented as a function of stiffness and a basis for practical realization. Target of motion is parameterized by control stiffness which could be variable according to the task reference. The system robustness of motion always requires very high stiffness in the controller.
Industrial
The branch of industrial engineer includes the design of machinery, assembly and process lines of various manufacturing industries. This branch can be said somewhat similar to automation and robotics. Mechatronics engineers who works as industrial engineers design and develop infrastructure of a manufacturing plant. Also it can be said that they are architect of machines. One can work as an industrial designer to design the industrial layout and plan for setting up of a manufacturing industry or as an industrial technician to lookover the technical requirements and repairing of the particular factory.
Robotics
Robotics is one of the newest emerging subfield of mechatronics. It is the study of robots and how they are manufactured and operated. Since 2000, this branch of mechatronics is attracting a number of aspirants. Robotics is interrelated with automation because here also not much human intervention is required. In a large number of factories, especially in automobile factories, robots are found in assembly lines, where they perform the job of drilling, installation and fitting. Programming skills are necessary for specialization in robotics. Knowledge of programming language—ROBOTC—is important for functioning robots. An industrial robot is a prime example of a mechatronics system; it includes aspects of electronics, mechanics and computing to do its day-to-day jobs.
Computer
The Internet of things (IoT) is the inter-networking of physical devices, embedded with electronics, software, sensors, actuators, and network connectivity which enable these objects to collect and exchange data. IoT and mechatronics are complementary. Many of the smart components associated with the Internet of Things will be essentially mechatronic. The development of the IoT is forcing mechatronics engineers, designers, practitioners and educators to research the ways in which mechatronic systems and components are perceived, designed and manufactured. This allows them to face up to new issues such as data security, machine ethics and the human-machine interface.
Knowledge of programming is very important. A mechatronics engineer has to do programming in different levels – for example, PLC programming, drone programming, hardware programming, CNC programming, etc. Due to combination of electronics engineering, soft skills from computer side is important. Important programming languages for mechatronics engineer to learn are Java, Python, Rust, C++ and C programming language.
| Technology | Disciplines | null |
304684 | https://en.wikipedia.org/wiki/Agaricus%20bisporus | Agaricus bisporus | Agaricus bisporus, commonly known as the cultivated mushroom, is a basidiomycete mushroom native to grasslands in Eurasia and North America. It is cultivated in more than 70 countries and is one of the most commonly and widely consumed mushrooms in the world. It has two color states while immature – white and brown – both of which have various names, with additional names for the mature state, such as chestnut, portobello, portabellini, button and champignon de Paris.
A. bisporus has some deadly poisonous lookalikes in the wild, such as Entoloma sinuatum.
Description
The pileus or cap of the original wild species is a pale grey-brown, with broad, flat scales on a paler background and fading toward the margins. It is first hemispherical before flattening out with maturity, and measures
in diameter. The narrow, crowded gills are free and initially pink, then red-brown, and finally a dark brown with a whitish edge from the cheilocystidia. The cylindrical stipe is up to tall by wide and bears a thick and narrow ring, which may be streaked on the upper side. The firm flesh is white, although it stains a pale pinkish-red on bruising. The spore print is dark brown. The spores are oval to round and measure approximately 4.5–5.5 μm × 5–7.5 μm, and the basidia usually two-spored, although two four-spored varieties have been described from the Mojave Desert and the Mediterranean, with predominantly heterothallic and homothallic lifestyles, respectively.
Similar species
The common mushroom could be confused with young specimens of a group of lethal mushrooms in the Amanita genus referred to as destroying angels, but the latter may be distinguished by their volva or cup at the base of the mushroom and pure white gills (as opposed to pinkish or brown of A. bisporus).
A more common and less dangerous mistake is to confuse Agaricus bisporus with A. xanthodermus, an inedible mushroom found worldwide in grassy areas. A. xanthodermus has an odor reminiscent of phenol; its flesh turns yellow when bruised. This fungus causes nausea and vomiting in some people.
The poisonous European species Entoloma sinuatum has a passing resemblance but has yellowish gills, turning pink, and lacks a ring.
Taxonomy
The common mushroom has a complicated taxonomic history. It was first described by English botanist Mordecai Cubitt Cooke in his 1871 Handbook of British Fungi, as a variety (var. hortensis) of Agaricus campestris. Danish mycologist Jakob Emanuel Lange later reviewed a cultivar specimen, and dubbed it Psalliota hortensis var. bispora in 1926. In 1938, it was promoted to species status and renamed Psalliota bispora. Emil Imbach (1897–1970) imparted the current scientific name of the species, Agaricus bisporus after the genus Psalliota was renamed to Agaricus in 1946. The specific epithet bispora distinguishes the two-spored basidia from four-spored varieties.
Names
When immature and , this mushroom may be known as:
common mushroom
white mushroom
button mushroom
cultivated mushroom
table mushroom
champignon (French for mushroom) de Paris
When immature and , it may be known variously as:
Swiss brown mushroom
Roman brown mushroom
Italian brown mushroom
cremini (also crimini) mushroom
chestnut mushroom (not to be confused with Pholiota adiposa)
baby bella
When marketed in its mature state, the mushroom is brown with a cap measuring . This form is commonly sold under the names portobello, portabella, or portobella. The etymology is disputed.
Distribution and habitat
This mushroom is commonly found worldwide in fields and grassy areas following rain, from late spring to autumn, especially in association with manure.
Cultivation
Production
In 2022, world production of mushrooms (including truffles) was 48 million tonnes, led by China with 94% of the total (table). Japan and the United States were secondary producers.
History
The earliest scientific description of the commercial cultivation of A. bisporus was made by French botanist Joseph Pitton de Tournefort in 1707. French agriculturist Olivier de Serres noted that transplanting mushroom mycelia would lead to the propagation of more mushrooms.
Originally, cultivation was unreliable as mushroom growers would watch for good flushes of mushrooms in fields before digging up the mycelium and replanting them in beds of composted manure or inoculating 'bricks' of compressed litter, loam, and manure. Spawn collected this way contained pathogens, and crops would be infected or not grow. In 1893, sterilized, or pure culture, spawn was discovered and produced by the Pasteur Institute in Paris for cultivation on composted horse manure.
Modern commercial varieties of the common agaricus mushroom were originally light brown. The white mushroom was discovered in 1925 growing among a bed of brown mushrooms at the Keystone Mushroom Farm in Coatesville, Pennsylvania. Louis Ferdinand Lambert, the farm's owner and a mycologist by training, brought the white mushroom back to his laboratory. As with the reception of white bread, it was seen as a more attractive food item and became grown and distributed. Similar to the commercial development history of the navel orange and Red Delicious apple, cultures were grown from the mutant individuals. Most cream-colored store mushrooms marketed today are products of this 1925 chance natural mutation.
A. bisporus is cultivated in at least seventy countries worldwide.
Nutrition
In a 100-gram serving, raw white mushrooms provide of food energy and are an excellent source (20% or more of the Daily Value, DV) of the B vitamins riboflavin, niacin, and pantothenic acid (table). Fresh mushrooms are also a good source (10–19% DV) of the dietary minerals phosphorus and potassium (table).
While fresh A. bisporus only contains 0.2 micrograms (8 IU) of vitamin D per 100 g, the ergocalciferol (D2) content increases substantially to 11.2 micrograms (446 IU) after exposure to UV light.
A. Bisporus contains 0.4 g/kg fresh of agaritine, a mycotoxin.
Gallery
| Biology and health sciences | Edible fungi | Plants |
304875 | https://en.wikipedia.org/wiki/Blogger%20%28service%29 | Blogger (service) | Blogger is an American online content management system founded in 1999 that enables its users to write blogs with time-stamped entries. Pyra Labs developed it before being acquired by Google in 2003. Google hosts the blogs, which can be accessed through a subdomain of blogspot.com. Blogs can also be accessed from a user-owned custom domain (such as www.example.com) by using DNS facilities to direct a domain to Google's servers. A user can have up to 100 blogs or websites per account.
Blogger enabled users to publish blogs and websites to their own web hosting server via FTP until May 1, 2010. All such blogs and websites had to be redirected to a blogspot.com subdomain or point their own domain to Google's servers via DNS.
History
Pyra Labs launched Blogger on August 23, 1999. It is credited with popularizing the format as one of the first dedicated blog-publishing tools. Pyra Labs was purchased by Google in February 2003 for an undisclosed amount. Premium features, which Pyra had actually offered for a fee, were made free as a result of the takeover. Evan Williams, a co-founder of Pyra Labs, left Google in October 2004. Picasa was acquired by Google in 2004, and Picasa and its photo-sharing service Hello were incorporated into Blogger, enabling users to upload images to their blogs.
Blogger underwent a major redesign on May 9, 2004, which included web standards-compliant templates, individual archive pages for posts, comments, and email posting. Blogger's new version, codenamed "Invader," was released in beta alongside the gold update on August 14, 2006. Users were moved to Google servers, and new features such as interface language in French, Italian, German, and Spanish were added. In December 2006, this new version of Blogger was taken out of beta. By May 2007, Blogger had completely moved over to Google-operated servers. Blogger was ranked 16 on the list of top 50 domains in terms of number of unique visitors in 2007.
On February 24, 2015, Blogger announced that as of late March it would no longer allow its users to post sexually explicit content, unless the nudity offers "substantial public benefit," for example in "artistic, educational, documentary, or scientific contexts." On February 28, 2015, accounting for severe backlash from long-term bloggers, Blogger reversed its decision on banning sexual content, going back to the previous policy that allowed explicit images and videos if the blog was marked as "adult".
Redesign
As part of the Blogger redesign in 2006, all blogs associated with a user's Google Account were migrated to Google servers. Blogger claims that the service is now more reliable because of the quality of the servers.
Along with the migration to Google servers, several new features were introduced, including label organization, a drag-and-drop template editing interface, reading permissions (to create private blogs) and new Web feed options. Furthermore, blogs are updated dynamically, as opposed to rewriting HTML files.
In a version of the service called Blogger in Draft, new features are tested before being released to all users. New features are discussed in the service's official blog.
In September 2009, Google introduced new features into Blogger as part of its tenth-anniversary celebration. The features included a new interface for post editing, improved image handling, Raw HTML Conversion, and other Google Docs-based implementations, including:
Adding location to posts via geotagging.
Post time-stamping at publication, not at original creation.
Vertical re-sizing of the post editor. The size is saved in a per-user, per-blog preference.
Link editing in compose mode.
Full Safari 3 support and fidelity on both Windows and macOS.
New Preview dialog that shows posts in a width and font size approximating what is seen in the published view.
Placeholder image for tags so that embeds are movable in compose mode.
New toolbar with Google aesthetics, faster loading time, and "undo" and "redo" buttons, also added the full justification button, a strike-through button, and an expanded color palette.
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304942 | https://en.wikipedia.org/wiki/Heart%20rate | Heart rate | Heart rate is the frequency of the heartbeat measured by the number of contractions of the heart per minute (beats per minute, or bpm). The heart rate varies according to the body's physical needs, including the need to absorb oxygen and excrete carbon dioxide. It is also modulated by numerous factors, including (but not limited to) genetics, physical fitness, stress or psychological status, diet, drugs, hormonal status, environment, and disease/illness, as well as the interaction between these factors. It is usually equal or close to the pulse rate measured at any peripheral point.
The American Heart Association states the normal resting adult human heart rate is 60–100 bpm. An ultra-trained athlete would have a resting heart rate of 37–38 bpm. Tachycardia is a high heart rate, defined as above 100 bpm at rest. Bradycardia is a low heart rate, defined as below 60 bpm at rest. When a human sleeps, a heartbeat with rates around 40–50 bpm is common and considered normal. When the heart is not beating in a regular pattern, this is referred to as an arrhythmia. Abnormalities of heart rate sometimes indicate disease.
Physiology
While heart rhythm is regulated entirely by the sinoatrial node under normal conditions, heart rate is regulated by sympathetic and parasympathetic input to the sinoatrial node. The accelerans nerve provides sympathetic input to the heart by releasing norepinephrine onto the cells of the sinoatrial node (SA node), and the vagus nerve provides parasympathetic input to the heart by releasing acetylcholine onto sinoatrial node cells. Therefore, stimulation of the accelerans nerve increases heart rate, while stimulation of the vagus nerve decreases it.
As water and blood are incompressible fluids, one of the physiological ways to deliver more blood to an organ is to increase heart rate. Normal resting heart rates range from 60 to 100 bpm. Bradycardia is defined as a resting heart rate below 60 bpm. However, heart rates from 50 to 60 bpm are common among healthy people and do not necessarily require special attention. Tachycardia is defined as a resting heart rate above 100 bpm, though persistent rest rates between 80 and 100 bpm, mainly if they are present during sleep, may be signs of hyperthyroidism or anemia (see below).
Central nervous system stimulants such as substituted amphetamines increase heart rate.
Central nervous system depressants or sedatives decrease the heart rate (apart from some particularly strange ones with equally strange effects, such as ketamine which can cause – amongst many other things – stimulant-like effects such as tachycardia).
There are many ways in which the heart rate speeds up or slows down. Most involve stimulant-like endorphins and hormones being released in the brain, some of which are those that are 'forced'/'enticed' out by the ingestion and processing of drugs such as cocaine or atropine.
This section discusses target heart rates for healthy persons, which would be inappropriately high for most persons with coronary artery disease.
Influences from the central nervous system
Cardiovascular centres
The heart rate is rhythmically generated by the sinoatrial node. It is also influenced by central factors through sympathetic and parasympathetic nerves. Nervous influence over the heart rate is centralized within the two paired cardiovascular centres of the medulla oblongata. The cardioaccelerator regions stimulate activity via sympathetic stimulation of the cardioaccelerator nerves, and the cardioinhibitory centers decrease heart activity via parasympathetic stimulation as one component of the vagus nerve. During rest, both centers provide slight stimulation to the heart, contributing to autonomic tone. This is a similar concept to tone in skeletal muscles. Normally, vagal stimulation predominates as, left unregulated, the SA node would initiate a sinus rhythm of approximately 100 bpm.
Both sympathetic and parasympathetic stimuli flow through the paired cardiac plexus near the base of the heart. The cardioaccelerator center also sends additional fibers, forming the cardiac nerves via sympathetic ganglia (the cervical ganglia plus superior thoracic ganglia T1–T4) to both the SA and AV nodes, plus additional fibers to the atria and ventricles. The ventricles are more richly innervated by sympathetic fibers than parasympathetic fibers. Sympathetic stimulation causes the release of the neurotransmitter norepinephrine (also known as noradrenaline) at the neuromuscular junction of the cardiac nerves. This shortens the repolarization period, thus speeding the rate of depolarization and contraction, which results in an increased heartrate. It opens chemical or ligand-gated sodium and calcium ion channels, allowing an influx of positively charged ions.
Norepinephrine binds to the beta–1 receptor. High blood pressure medications are used to block these receptors and so reduce the heart rate.
Parasympathetic stimulation originates from the cardioinhibitory region of the brain with impulses traveling via the vagus nerve (cranial nerve X). The vagus nerve sends branches to both the SA and AV nodes, and to portions of both the atria and ventricles. Parasympathetic stimulation releases the neurotransmitter acetylcholine (ACh) at the neuromuscular junction. ACh slows HR by opening chemical- or ligand-gated potassium ion channels to slow the rate of spontaneous depolarization, which extends repolarization and increases the time before the next spontaneous depolarization occurs. Without any nervous stimulation, the SA node would establish a sinus rhythm of approximately 100 bpm. Since resting rates are considerably less than this, it becomes evident that parasympathetic stimulation normally slows HR. This is similar to an individual driving a car with one foot on the brake pedal. To speed up, one need merely remove one's foot from the brake and let the engine increase speed. In the case of the heart, decreasing parasympathetic stimulation decreases the release of ACh, which allows HR to increase up to approximately 100 bpm. Any increases beyond this rate would require sympathetic stimulation.
Input to the cardiovascular centres
The cardiovascular centre receive input from a series of visceral receptors with impulses traveling through visceral sensory fibers within the vagus and sympathetic nerves via the cardiac plexus. Among these receptors are various proprioreceptors, baroreceptors, and chemoreceptors, plus stimuli from the limbic system which normally enable the precise regulation of heart function, via cardiac reflexes. Increased physical activity results in increased rates of firing by various proprioreceptors located in muscles, joint capsules, and tendons. The cardiovascular centres monitor these increased rates of firing, suppressing parasympathetic stimulation or increasing sympathetic stimulation as needed in order to increase blood flow.
Similarly, baroreceptors are stretch receptors located in the aortic sinus, carotid bodies, the venae cavae, and other locations, including pulmonary vessels and the right side of the heart itself. Rates of firing from the baroreceptors represent blood pressure, level of physical activity, and the relative distribution of blood. The cardiac centers monitor baroreceptor firing to maintain cardiac homeostasis, a mechanism called the baroreceptor reflex. With increased pressure and stretch, the rate of baroreceptor firing increases, and the cardiac centers decrease sympathetic stimulation and increase parasympathetic stimulation. As pressure and stretch decrease, the rate of baroreceptor firing decreases, and the cardiac centers increase sympathetic stimulation and decrease parasympathetic stimulation.
There is a similar reflex, called the atrial reflex or Bainbridge reflex, associated with varying rates of blood flow to the atria. Increased venous return stretches the walls of the atria where specialized baroreceptors are located. However, as the atrial baroreceptors increase their rate of firing and as they stretch due to the increased blood pressure, the cardiac center responds by increasing sympathetic stimulation and inhibiting parasympathetic stimulation to increase HR. The opposite is also true.
Increased metabolic byproducts associated with increased activity, such as carbon dioxide, hydrogen ions, and lactic acid, plus falling oxygen levels, are detected by a suite of chemoreceptors innervated by the glossopharyngeal and vagus nerves. These chemoreceptors provide feedback to the cardiovascular centers about the need for increased or decreased blood flow, based on the relative levels of these substances.
The limbic system can also significantly impact HR related to emotional state. During periods of stress, it is not unusual to identify higher than normal HRs, often accompanied by a surge in the stress hormone cortisol. Individuals experiencing extreme anxiety may manifest panic attacks with symptoms that resemble those of heart attacks. These events are typically transient and treatable. Meditation techniques have been developed to ease anxiety and have been shown to lower HR effectively. Doing simple deep and slow breathing exercises with one's eyes closed can also significantly reduce this anxiety and HR.
Factors influencing heart rate
Using a combination of autorhythmicity and innervation, the cardiovascular center is able to provide relatively precise control over the heart rate, but other factors can impact on this. These include hormones, notably epinephrine, norepinephrine, and thyroid hormones; levels of various ions including calcium, potassium, and sodium; body temperature; hypoxia; and pH balance.
Epinephrine and norepinephrine
The catecholamines, epinephrine and norepinephrine, secreted by the adrenal medulla form one component of the extended fight-or-flight mechanism. The other component is sympathetic stimulation. Epinephrine and norepinephrine have similar effects: binding to the beta-1 adrenergic receptors, and opening sodium and calcium ion chemical- or ligand-gated channels. The rate of depolarization is increased by this additional influx of positively charged ions, so the threshold is reached more quickly and the period of repolarization is shortened. However, massive releases of these hormones coupled with sympathetic stimulation may actually lead to arrhythmias. There is no parasympathetic stimulation to the adrenal medulla.
Thyroid hormones
In general, increased levels of the thyroid hormones (thyroxine(T4) and triiodothyronine (T3)), increase the heart rate; excessive levels can trigger tachycardia. The impact of thyroid hormones is typically of a much longer duration than that of the catecholamines. The physiologically active form of triiodothyronine, has been shown to directly enter cardiomyocytes and alter activity at the level of the genome. It also impacts the beta-adrenergic response similar to epinephrine and norepinephrine.
Calcium
Calcium ion levels have a great impact on heart rate and myocardial contractility: increased calcium levels cause an increase in both. High levels of calcium ions result in hypercalcemia and excessive levels can induce cardiac arrest. Drugs known as calcium channel blockers slow HR by binding to these channels and blocking or slowing the inward movement of calcium ions.
Caffeine and nicotine
Caffeine and nicotine are both stimulants of the nervous system and of the cardiac centres causing an increased heart rate. Caffeine works by increasing the rates of depolarization at the SA node, whereas nicotine stimulates the activity of the sympathetic neurons that deliver impulses to the heart.
Effects of stress
Both surprise and stress induce physiological response: elevate heart rate substantially. In a study conducted on 8 female and male student actors ages 18 to 25, their reaction to an unforeseen occurrence (the cause of stress) during a performance was observed in terms of heart rate. In the data collected, there was a noticeable trend between the location of actors (onstage and offstage) and their elevation in heart rate in response to stress; the actors present offstage reacted to the stressor immediately, demonstrated by their immediate elevation in heart rate the minute the unexpected event occurred, but the actors present onstage at the time of the stressor reacted in the following 5 minute period (demonstrated by their increasingly elevated heart rate). This trend regarding stress and heart rate is supported by previous studies; negative emotion/stimulus has a prolonged effect on heart rate in individuals who are directly impacted.
In regard to the characters present onstage, a reduced startle response has been associated with a passive defense, and the diminished initial heart rate response has been predicted to have a greater tendency to dissociation. Current evidence suggests that heart rate variability can be used as an accurate measure of psychological stress and may be used for an objective measurement of psychological stress.
Factors decreasing heart rate
The heart rate can be slowed by altered sodium and potassium levels, hypoxia, acidosis, alkalosis, and hypothermia. The relationship between electrolytes and HR is complex, but maintaining electrolyte balance is critical to the normal wave of depolarization. Of the two ions, potassium has the greater clinical significance. Initially, both hyponatremia (low sodium levels) and hypernatremia (high sodium levels) may lead to tachycardia. Severely high hypernatremia may lead to fibrillation, which may cause cardiac output to cease. Severe hyponatremia leads to both bradycardia and other arrhythmias. Hypokalemia (low potassium levels) also leads to arrhythmias, whereas hyperkalemia (high potassium levels) causes the heart to become weak and flaccid, and ultimately to fail.
Heart muscle relies exclusively on aerobic metabolism for energy. Severe myocardial infarction (commonly called a heart attack) can lead to a decreasing heart rate, since metabolic reactions fueling heart contraction are restricted.
Acidosis is a condition in which excess hydrogen ions are present, and the patient's blood expresses a low pH value. Alkalosis is a condition in which there are too few hydrogen ions, and the patient's blood has an elevated pH. Normal blood pH falls in the range of 7.35–7.45, so a number lower than this range represents acidosis and a higher number represents alkalosis. Enzymes, being the regulators or catalysts of virtually all biochemical reactions – are sensitive to pH and will change shape slightly with values outside their normal range. These variations in pH and accompanying slight physical changes to the active site on the enzyme decrease the rate of formation of the enzyme-substrate complex, subsequently decreasing the rate of many enzymatic reactions, which can have complex effects on HR. Severe changes in pH will lead to denaturation of the enzyme.
The last variable is body temperature. Elevated body temperature is called hyperthermia, and suppressed body temperature is called hypothermia. Slight hyperthermia results in increasing HR and strength of contraction. Hypothermia slows the rate and strength of heart contractions. This distinct slowing of the heart is one component of the larger diving reflex that diverts blood to essential organs while submerged. If sufficiently chilled, the heart will stop beating, a technique that may be employed during open heart surgery. In this case, the patient's blood is normally diverted to an artificial heart-lung machine to maintain the body's blood supply and gas exchange until the surgery is complete, and sinus rhythm can be restored. Excessive hyperthermia and hypothermia will both result in death, as enzymes drive the body systems to cease normal function, beginning with the central nervous system.
Physiological control over heart rate
A study shows that bottlenose dolphins can learn – apparently via instrumental conditioning – to rapidly and selectively slow down their heart rate during diving for conserving oxygen depending on external signals. In humans regulating heart rate by methods such as listening to music, meditation or a vagal maneuver takes longer and only lowers the rate to a much smaller extent.
In different circumstances
Heart rate is not a stable value and it increases or decreases in response to the body's need in a way to maintain an equilibrium (basal metabolic rate) between requirement and delivery of oxygen and nutrients. The normal SA node firing rate is affected by autonomic nervous system activity: sympathetic stimulation increases and parasympathetic stimulation decreases the firing rate.
Resting heart rate
Normal pulse rates at rest, in beats per minute (BPM):
The basal or resting heart rate (HRrest) is defined as the heart rate when a person is awake, in a neutrally temperate environment, and has not been subject to any recent exertion or stimulation, such as stress or surprise. The normal resting heart rate is based on the at-rest firing rate of the heart's sinoatrial node, where the faster pacemaker cells driving the self-generated rhythmic firing and responsible for the heart's autorhythmicity are located.
In one study 98% of cardiologists suggested that as a desirable target range, 50 to 90 beats per minute is more appropriate than 60 to 100. The available evidence indicates that the normal range for resting heart rate is 50–90 beats per minute (bpm). In a study of over 35,000 American men and women over age 40 during the 1999–2008 period, 71 bpm was the average for men, and 73 bpm was the average for women.
Resting heart rate is often correlated with mortality. In the Copenhagen City Heart Study a heart rate of 65 bpm rather than 80 bpm was associated with 4.6 years longer life expectancy in men and 3.6 years in women. Other studies have shown all-cause mortality is increased by 1.22 (hazard ratio) when heart rate exceeds 90 beats per minute. ECG of 46,129 individuals with low risk for cardiovascular disease revealed that 96% had resting heart rates ranging from 48 to 98 beats per minute. The mortality rate of patients with myocardial infarction increased from 15% to 41% if their admission heart rate was greater than 90 beats per minute. For endurance athletes at the elite level, it is not unusual to have a resting heart rate between 33 and 50 bpm.
Maximum heart rate
The maximum heart rate (HRmax) is the age-related highest number of beats per minute of the heart when reaching a point of exhaustion without severe problems through exercise stress.
In general it is loosely estimated as 220 minus one's age.
It generally decreases with age. Since HRmax varies by individual, the most accurate way of measuring any single person's HRmax is via a cardiac stress test. In this test, a person is subjected to controlled physiologic stress (generally by treadmill or bicycle ergometer) while being monitored by an electrocardiogram (ECG). The intensity of exercise is periodically increased until certain changes in heart function are detected on the ECG monitor, at which point the subject is directed to stop. Typical duration of the test ranges ten to twenty minutes. Adults who are beginning a new exercise regimen are often advised to perform this test only in the presence of medical staff due to risks associated with high heart rates.
The theoretical maximum heart rate of a human is 300 bpm; however, there have been multiple cases where this theoretical upper limit has been exceeded. The fastest human ventricular conduction rate recorded to this day is a conducted tachyarrhythmia with ventricular rate of 600 beats per minute, which is comparable to the heart rate of a mouse.
For general purposes, a number of formulas are used to estimate HRmax. However, these predictive formulas have been criticized as inaccurate because they only produce generalized population-averages and may deviate significantly from the actual value. (See § Limitations.)
Haskell & Fox (1970)
Notwithstanding later research, the most widely cited formula for HRmax is still:
HRmax = 220 − age
Although attributed to various sources, it is widely thought to have been devised in 1970 by Dr. William Haskell and Dr. Samuel Fox. They did not develop this formula from original research, but rather by plotting data from approximately 11 references consisting of published research or unpublished scientific compilations. It gained widespread use through being used by Polar Electro in its heart rate monitors, which Dr. Haskell has "laughed about", as the formula "was never supposed to be an absolute guide to rule people's training."
While this formula is commonly used (and easy to remember and calculate), research has consistently found that it is subject to bias, particularly in older adults. Compared to the age-specific average HRmax, the Haskell and Fox formula overestimates HRmax in young adults, agrees with it at age 40, and underestimates HRmax in older adults. For example, in one study, the average HRmax at age 76 was about 10bpm higher than the Haskell and Fox equation. Consequently, the formula cannot be recommended for use in exercise physiology and related fields.
Other formulas
HRmax is strongly correlated to age, and most formulas are solely based on this. Studies have been mixed on the effect of gender, with some finding that gender is statistically significant, although small when considering overall equation error, while others finding negligible effect. The inclusion of physical activity status, maximal oxygen uptake, smoking, body mass index, body weight, or resting heart rate did not significantly improve accuracy. Nonlinear models are slightly more accurate predictors of average age-specific HRmax, particularly above 60 years of age, but are harder to apply, and provide statistically negligible improvement over linear models. The Wingate formula is the most recent, had the largest data set, and performed best on a fresh data set when compared with other formulas, although it had only a small amount of data for ages 60 and older so those estimates should be viewed with caution. In addition, most formulas are developed for adults and are not applicable to children and adolescents.
Limitations
Maximum heart rates vary significantly between individuals. Age explains only about half of HRmax variance. For a given age, the standard deviation of HRmax from the age-specific population mean is about 12bpm, and a 95% interval for the prediction error is about 24bpm. For example, Dr. Fritz Hagerman observed that the maximum heart rates of men in their 20s on Olympic rowing teams vary from 160 to 220. Such a variation would equate to an age range of -16 to 68 using the Wingate formula. The formulas are quite accurate at predicting the average heart rate of a group of similarly-aged individuals, but relatively poor for a given individual.
Robergs and Landwehr opine that for VO2 max, prediction errors in HRmax need to be less than ±3 bpm. No current formula meets this accuracy. For prescribing exercise training heart rate ranges, the errors in the more accurate formulas may be acceptable, but again it is likely that, for a significant fraction of the population, current equations used to estimate HRmax are not accurate enough. Froelicher and Myers describe maximum heart formulas as "largely useless". Measurement via a maximal test is preferable whenever possible, which can be as accurate as ±2bpm.
Heart rate reserve
Heart rate reserve (HRreserve) is the difference between a person's measured or predicted maximum heart rate and resting heart rate. Some methods of measurement of exercise intensity measure percentage of heart rate reserve. Additionally, as a person increases their cardiovascular fitness, their HRrest will drop, and the heart rate reserve will increase. Percentage of HRreserve is statistically indistinguishable from percentage of VO2 reserve.
HRreserve = HRmax − HRrest
This is often used to gauge exercise intensity (first used in 1957 by Karvonen).
Karvonen's study findings have been questioned, due to the following:
The study did not use VO2 data to develop the equation.
Only six subjects were used.
Karvonen incorrectly reported that the percentages of HRreserve and VO2 max correspond to each other, but newer evidence shows that it correlated much better with VO2 reserve as described above.
Target heart rate
For healthy people, the Target Heart Rate (THR) or Training Heart Rate Range (THRR) is a desired range of heart rate reached during aerobic exercise which enables one's heart and lungs to receive the most benefit from a workout. This theoretical range varies based mostly on age; however, a person's physical condition, sex, and previous training also are used in the calculation.
By percent, Fox–Haskell-based
The THR can be calculated as a range of 65–85% intensity, with intensity defined simply as percentage of HRmax. However, it is crucial to derive an accurate HRmax to ensure these calculations are meaningful.
Example for someone with a HRmax of 180 (age 40, estimating HRmax As 220 − age):
65% Intensity: (220 − (age = 40)) × 0.65 → 117 bpm
85% Intensity: (220 − (age = 40)) × 0.85 → 154 bpm
Karvonen method
The Karvonen method factors in resting heart rate (HRrest) to calculate target heart rate (THR), using a range of 50–85% intensity:
THR = ((HRmax − HRrest) × % intensity) + HRrest
Equivalently,
THR = (HRreserve × % intensity) + HRrest
Example for someone with a HRmax of 180 and a HRrest of 70 (and therefore a HRreserve of 110):
50% Intensity: ((180 − 70) × 0.50) + 70 = 125 bpm
85% Intensity: ((180 − 70) × 0.85) + 70 = 163 bpm
Zoladz method
An alternative to the Karvonen method is the Zoladz method, which is used to test an athlete's capabilities at specific heart rates. These are not intended to be used as exercise zones, although they are often used as such. The Zoladz test zones are derived by subtracting values from HRmax:
THR = HRmax − Adjuster ± 5 bpm
Zone 1 Adjuster = 50 bpm
Zone 2 Adjuster = 40 bpm
Zone 3 Adjuster = 30 bpm
Zone 4 Adjuster = 20 bpm
Zone 5 Adjuster = 10 bpm
Example for someone with a HRmax of 180:
Zone 1 (easy exercise): 180 − 50 ± 5 → 125 − 135 bpm
Zone 4 (tough exercise): 180 − 20 ± 5 → 155 − 165 bpm
Heart rate recovery
Heart rate recovery (HRR) is the reduction in heart rate at peak exercise and the rate as measured after a cool-down period of fixed duration. A greater reduction in heart rate after exercise during the reference period is associated with a higher level of cardiac fitness.
Heart rates assessed during treadmill stress test that do not drop by more than 12 bpm one minute after stopping exercise (if cool-down period after exercise) or by more than 18 bpm one minute after stopping exercise (if no cool-down period and supine position as soon as possible) are associated with an increased risk of death. People with an abnormal HRR defined as a decrease of 42 beats per minutes or less at two minutes post-exercise had a mortality rate 2.5 times greater than patients with a normal recovery. Another study reported a four-fold increase in mortality in subjects with an abnormal HRR defined as ≤12 bpm reduction one minute after the cessation of exercise. A study reported that a HRR of ≤22 bpm after two minutes "best identified high-risk patients". They also found that while HRR had significant prognostic value it had no diagnostic value.
Development
The human heart beats more than 2.8 billion times in an average lifetime.
The heartbeat of a human embryo begins at approximately 21 days after conception, or five weeks after the last normal menstrual period (LMP), which is the date normally used to date pregnancy in the medical community. The electrical depolarizations that trigger cardiac myocytes to contract arise spontaneously within the myocyte itself. The heartbeat is initiated in the pacemaker regions and spreads to the rest of the heart through a conduction pathway. Pacemaker cells develop in the primitive atrium and the sinus venosus to form the sinoatrial node and the atrioventricular node respectively. Conductive cells develop the bundle of His and carry the depolarization into the lower heart.
The human heart begins beating at a rate near the mother's, about 75–80 beats per minute (bpm). The embryonic heart rate then accelerates linearly for the first month of beating, peaking at 165–185 bpm during the early 7th week, (early 9th week after the LMP). This acceleration is approximately 3.3 bpm per day, or about 10 bpm every three days, an increase of 100 bpm in the first month.
After peaking at about 9.2 weeks after the LMP, it decelerates to about 150 bpm (+/-25 bpm) during the 15th week after the LMP. After the 15th week the deceleration slows reaching an average rate of about 145 (+/-25 bpm) bpm at term. The regression formula which describes this acceleration before the embryo reaches 25 mm in crown-rump length or 9.2 LMP weeks is:
Clinical significance
Manual measurement
Heart rate is measured by finding the pulse of the heart. This pulse rate can be found at any point on the body where the artery's pulsation is transmitted to the surface by pressuring it with the index and middle fingers; often it is compressed against an underlying structure like bone. The thumb should not be used for measuring another person's heart rate, as its strong pulse may interfere with the correct perception of the target pulse.
The radial artery is the easiest to use to check the heart rate. However, in emergency situations the most reliable arteries to measure heart rate are carotid arteries. This is important mainly in patients with atrial fibrillation, in whom heart beats are irregular and stroke volume is largely different from one beat to another. In those beats following a shorter diastolic interval left ventricle does not fill properly, stroke volume is lower and pulse wave is not strong enough to be detected by palpation on a distal artery like the radial artery. It can be detected, however, by doppler.
Possible points for measuring the heart rate are:
The ventral aspect of the wrist on the side of the thumb (radial artery).
The ulnar artery.
The inside of the elbow, or under the biceps muscle (brachial artery).
The groin (femoral artery).
Behind the medial malleolus on the feet (posterior tibial artery).
Middle of dorsum of the foot (dorsalis pedis).
Behind the knee (popliteal artery).
Over the abdomen (abdominal aorta).
The chest (apex of the heart), which can be felt with one's hand or fingers. It is also possible to auscultate the heart using a stethoscope.
In the neck, lateral of the larynx (carotid artery)
The temple (superficial temporal artery).
The lateral edge of the mandible (facial artery).
The side of the head near the ear (posterior auricular artery).
Electronic measurement
A more precise method of determining heart rate involves the use of an electrocardiograph, or ECG (also abbreviated EKG). An ECG generates a pattern based on electrical activity of the heart, which closely follows heart function. Continuous ECG monitoring is routinely done in many clinical settings, especially in critical care medicine. On the ECG, instantaneous heart rate is calculated using the R wave-to-R wave (RR) interval and multiplying/dividing in order to derive heart rate in heartbeats/min. Multiple methods exist:
HR = 1000 · 60/(RR interval in milliseconds)
HR = 60/(RR interval in seconds)
HR = 300/number of "large" squares between successive R waves.
HR= 1,500 number of large blocks
Heart rate monitors allow measurements to be taken continuously and can be used during exercise when manual measurement would be difficult or impossible (such as when the hands are being used). Various commercial heart rate monitors are also available. Some monitors, used during sport, consist of a chest strap with electrodes. The signal is transmitted to a wrist receiver for display.
Alternative methods of measurement include seismocardiography.
Optical measurements
Pulse oximetry of the finger and laser Doppler imaging of the eye fundus are often used in the clinics. Those techniques can assess the heart rate by measuring the delay between pulses.
Tachycardia
Tachycardia is a resting heart rate more than 100 beats per minute. This number can vary as smaller people and children have faster heart rates than average adults.
Physiological conditions where tachycardia occurs:
Pregnancy
Emotional conditions such as anxiety or stress.
Exercise
Pathological conditions where tachycardia occurs:
Sepsis
Fever
Anemia
Hypoxia
Hyperthyroidism
Hypersecretion of catecholamines
Cardiomyopathy
Valvular heart diseases
Acute Radiation Syndrome
Dehydration
Metabolic myopathies (At rest, tachycardia is commonly seen in fatty acid oxidation disorders. An inappropriate rapid heart rate response to exercise is seen in muscle glycogenoses and mitochondrial myopathies, where the tachycardia is faster than would be expected during exercise).
Bradycardia
Bradycardia was defined as a heart rate less than 60 beats per minute when textbooks asserted that the normal range for heart rates was 60–100 bpm. The normal range has since been revised in textbooks to 50–90 bpm for a human at total rest. Setting a lower threshold for bradycardia prevents misclassification of fit individuals as having a pathologic heart rate. The normal heart rate number can vary as children and adolescents tend to have faster heart rates than average adults. Bradycardia may be associated with medical conditions such as hypothyroidism, heart disease, or inflammatory disease. At rest, although tachycardia is more commonly seen in fatty acid oxidation disorders, more rarely acute bradycardia can occur.
Trained athletes tend to have slow resting heart rates, and resting bradycardia in athletes should not be considered abnormal if the individual has no symptoms associated with it. For example, Miguel Indurain, a Spanish cyclist and five time Tour de France winner, had a resting heart rate of 28 beats per minute, one of the lowest ever recorded in a healthy human. Daniel Green achieved the world record for the slowest heartbeat in a healthy human with a heart rate of just 26 bpm in 2014.
Arrhythmia
Arrhythmias are abnormalities of the heart rate and rhythm (sometimes felt as palpitations). They can be divided into two broad categories: fast and slow heart rates. Some cause few or minimal symptoms. Others produce more serious symptoms of lightheadedness, dizziness and fainting.
Hypertension
Elevated heart rate is a powerful predictor of morbidity and mortality in patients with hypertension. Atherosclerosis and dysautonomia are major contributors to the pathogenesis.
Correlation with cardiovascular mortality risk
A number of investigations indicate that faster resting heart rate has emerged as a new risk factor for mortality in homeothermic mammals, particularly cardiovascular mortality in human beings. High heart rate is associated with endothelial dysfunction and increased atheromatous plaque formation leading to atherosclerosis. Faster heart rate may accompany increased production of inflammation molecules and increased production of reactive oxygen species in cardiovascular system, in addition to increased mechanical stress to the heart. There is a correlation between increased resting rate and cardiovascular risk. This is not seen to be "using an allotment of heart beats" but rather an increased risk to the system from the increased rate.
An Australian-led international study of patients with cardiovascular disease has shown that heart beat rate is a key indicator for the risk of heart attack. The study, published in The Lancet (September 2008) studied 11,000 people, across 33 countries, who were being treated for heart problems. Those patients whose heart rate was above 70 beats per minute had significantly higher incidence of heart attacks, hospital admissions and the need for surgery. Higher heart rate is thought to be correlated with an increase in heart attack and about a 46 percent increase in hospitalizations for non-fatal or fatal heart attack.
Other studies have shown that a high resting heart rate is associated with an increase in cardiovascular and all-cause mortality in the general population and in patients with chronic diseases. A faster resting heart rate is associated with shorter life expectancy and is considered a strong risk factor for heart disease and heart failure, independent of level of physical fitness. Specifically, a resting heart rate above 65 beats per minute has been shown to have a strong independent effect on premature mortality; every 10 beats per minute increase in resting heart rate has been shown to be associated with a 10–20% increase in risk of death. In one study, men with no evidence of heart disease and a resting heart rate of more than 90 beats per minute had a five times higher risk of sudden cardiac death. Similarly, another study found that men with resting heart rates of over 90 beats per minute had an almost two-fold increase in risk for cardiovascular disease mortality; in women it was associated with a three-fold increase. In patients having heart rates of 70 beats/minute or above, each additional beat/minute was associated with increased rate of cardiovascular death and heart failure hospitalization.
Given these data, heart rate should be considered in the assessment of cardiovascular risk, even in apparently healthy individuals. Heart rate has many advantages as a clinical parameter: It is inexpensive and quick to measure and is easily understandable. Although the accepted limits of heart rate are between 60 and 100 beats per minute, this was based for convenience on the scale of the squares on electrocardiogram paper; a better definition of normal sinus heart rate may be between 50 and 90 beats per minute.
Standard textbooks of physiology and medicine mention that heart rate (HR) is readily calculated from the ECG as follows: HR = 1000*60/RR interval in milliseconds, HR = 60/RR interval in seconds, or HR = 300/number of large squares between successive R waves. In each case, the authors are actually referring to instantaneous HR, which is the number of times the heart would beat if successive RR intervals were constant.
Lifestyle and pharmacological regimens may be beneficial to those with high resting heart rates. Exercise is one possible measure to take when an individual's heart rate is higher than 80 beats per minute. Diet has also been found to be beneficial in lowering resting heart rate: In studies of resting heart rate and risk of death and cardiac complications on patients with type 2 diabetes, legumes were found to lower resting heart rate. This is thought to occur because in addition to the direct beneficial effects of legumes, they also displace animal proteins in the diet, which are higher in saturated fat and cholesterol. Another nutrient is omega-3 long chain polyunsaturated fatty acids (omega-3 fatty acid or LC-PUFA). In a meta-analysis with a total of 51 randomized controlled trials (RCTs) involving 3,000 participants, the supplement mildly but significantly reduced heart rate (-2.23 bpm; 95% CI: -3.07, -1.40 bpm). When docosahexaenoic acid (DHA) and eicosapentaenoic acid (EPA) were compared, modest heart rate reduction was observed in trials that supplemented with DHA (-2.47 bpm; 95% CI: -3.47, -1.46 bpm), but not in those received EPA.
A very slow heart rate (bradycardia) may be associated with heart block. It may also arise from autonomous nervous system impairment.
| Biology and health sciences | Symptoms and signs | Health |
305136 | https://en.wikipedia.org/wiki/Visual%20system | Visual system | The visual system is the physiological basis of visual perception (the ability to detect and process light). The system detects, transduces and interprets information concerning light within the visible range to construct an image and build a mental model of the surrounding environment. The visual system is associated with the eye and functionally divided into the optical system (including cornea and lens) and the neural system (including the retina and visual cortex).
The visual system performs a number of complex tasks based on the image forming functionality of the eye, including the formation of monocular images, the neural mechanisms underlying stereopsis and assessment of distances to (depth perception) and between objects, motion perception, pattern recognition, accurate motor coordination under visual guidance, and colour vision. Together, these facilitate higher order tasks, such as object identification. The neuropsychological side of visual information processing is known as visual perception, an abnormality of which is called visual impairment, and a complete absence of which is called blindness. The visual system also has several non-image forming visual functions, independent of visual perception, including the pupillary light reflex and circadian photoentrainment.
This article describes the human visual system, which is representative of mammalian vision, and to a lesser extent the vertebrate visual system.
System overview
Optical
Together, the cornea and lens refract light into a small image and shine it on the retina. The retina transduces this image into electrical pulses using rods and cones. The optic nerve then carries these pulses through the optic canal. Upon reaching the optic chiasm the nerve fibers decussate (left becomes right). The fibers then branch and terminate in three places.
Neural
Most of the optic nerve fibers end in the lateral geniculate nucleus (LGN). Before the LGN forwards the pulses to V1 of the visual cortex (primary) it gauges the range of objects and tags every major object with a velocity tag. These tags predict object movement.
The LGN also sends some fibers to V2 and V3.
V1 performs edge-detection to understand spatial organization (initially, 40 milliseconds in, focusing on even small spatial and color changes. Then, 100 milliseconds in, upon receiving the translated LGN, V2, and V3 info, also begins focusing on global organization). V1 also creates a bottom-up saliency map to guide attention or gaze shift.
V2 both forwards (direct and via pulvinar) pulses to V1 and receives them. Pulvinar is responsible for saccade and visual attention. V2 serves much the same function as V1, however, it also handles illusory contours, determining depth by comparing left and right pulses (2D images), and foreground distinguishment. V2 connects to V1 - V5.
V3 helps process 'global motion' (direction and speed) of objects. V3 connects to V1 (weak), V2, and the inferior temporal cortex.
V4 recognizes simple shapes, and gets input from V1 (strong), V2, V3, LGN, and pulvinar. V5's outputs include V4 and its surrounding area, and eye-movement motor cortices (frontal eye-field and lateral intraparietal area).
V5's functionality is similar to that of the other V's, however, it integrates local object motion into global motion on a complex level. V6 works in conjunction with V5 on motion analysis. V5 analyzes self-motion, whereas V6 analyzes motion of objects relative to the background. V6's primary input is V1, with V5 additions. V6 houses the topographical map for vision. V6 outputs to the region directly around it (V6A). V6A has direct connections to arm-moving cortices, including the premotor cortex.
The inferior temporal gyrus recognizes complex shapes, objects, and faces or, in conjunction with the hippocampus, creates new memories. The pretectal area is seven unique nuclei. Anterior, posterior and medial pretectal nuclei inhibit pain (indirectly), aid in REM, and aid the accommodation reflex, respectively. The Edinger-Westphal nucleus moderates pupil dilation and aids (since it provides parasympathetic fibers) in convergence of the eyes and lens adjustment. Nuclei of the optic tract are involved in smooth pursuit eye movement and the accommodation reflex, as well as REM.
The suprachiasmatic nucleus is the region of the hypothalamus that halts production of melatonin (indirectly) at first light.
Structure
The eye, especially the retina
The optic nerve
The optic chiasma
The optic tract
The lateral geniculate body
The optic radiation
The visual cortex
The visual association cortex.
These are components of the visual pathway, also called the optic pathway, that can be divided into anterior and posterior visual pathways. The anterior visual pathway refers to structures involved in vision before the lateral geniculate nucleus. The posterior visual pathway refers to structures after this point.
Eye
Light entering the eye is refracted as it passes through the cornea. It then passes through the pupil (controlled by the iris) and is further refracted by the lens. The cornea and lens act together as a compound lens to project an inverted image onto the retina.
Retina
The retina consists of many photoreceptor cells which contain particular protein molecules called opsins. In humans, two types of opsins are involved in conscious vision: rod opsins and cone opsins. (A third type, melanopsin in some retinal ganglion cells (RGC), part of the body clock mechanism, is probably not involved in conscious vision, as these RGC do not project to the lateral geniculate nucleus but to the pretectal olivary nucleus.) An opsin absorbs a photon (a particle of light) and transmits a signal to the cell through a signal transduction pathway, resulting in hyper-polarization of the photoreceptor.
Rods and cones differ in function. Rods are found primarily in the periphery of the retina and are used to see at low levels of light. Each human eye contains 120 million rods. Cones are found primarily in the center (or fovea) of the retina. There are three types of cones that differ in the wavelengths of light they absorb; they are usually called short or blue, middle or green, and long or red. Cones mediate day vision and can distinguish color and other features of the visual world at medium and high light levels. Cones are larger and much less numerous than rods (there are 6-7 million of them in each human eye).
In the retina, the photoreceptors synapse directly onto bipolar cells, which in turn synapse onto ganglion cells of the outermost layer, which then conduct action potentials to the brain. A significant amount of visual processing arises from the patterns of communication between neurons in the retina. About 130 million photo-receptors absorb light, yet roughly 1.2 million axons of ganglion cells transmit information from the retina to the brain. The processing in the retina includes the formation of center-surround receptive fields of bipolar and ganglion cells in the retina, as well as convergence and divergence from photoreceptor to bipolar cell. In addition, other neurons in the retina, particularly horizontal and amacrine cells, transmit information laterally (from a neuron in one layer to an adjacent neuron in the same layer), resulting in more complex receptive fields that can be either indifferent to color and sensitive to motion or sensitive to color and indifferent to motion.
Mechanism of generating visual signals
The retina adapts to change in light through the use of the rods. In the dark, the chromophore retinal has a bent shape called cis-retinal (referring to a cis conformation in one of the double bonds). When light interacts with the retinal, it changes conformation to a straight form called trans-retinal and breaks away from the opsin. This is called bleaching because the purified rhodopsin changes from violet to colorless in the light. At baseline in the dark, the rhodopsin absorbs no light and releases glutamate, which inhibits the bipolar cell. This inhibits the release of neurotransmitters from the bipolar cells to the ganglion cell. When there is light present, glutamate secretion ceases, thus no longer inhibiting the bipolar cell from releasing neurotransmitters to the ganglion cell and therefore an image can be detected.
The final result of all this processing is five different populations of ganglion cells that send visual (image-forming and non-image-forming) information to the brain:
M cells, with large center-surround receptive fields that are sensitive to depth, indifferent to color, and rapidly adapt to a stimulus;
P cells, with smaller center-surround receptive fields that are sensitive to color and shape;
K cells, with very large center-only receptive fields that are sensitive to color and indifferent to shape or depth;
another population that is intrinsically photosensitive; and
a final population that is used for eye movements.
A 2006 University of Pennsylvania study calculated the approximate bandwidth of human retinas to be about 8,960 kilobits per second, whereas guinea pig retinas transfer at about 875 kilobits.
In 2007 Zaidi and co-researchers on both sides of the Atlantic studying patients without rods and cones, discovered that the novel photoreceptive ganglion cell in humans also has a role in conscious and unconscious visual perception. The peak spectral sensitivity was 481 nm. This shows that there are two pathways for vision in the retina – one based on classic photoreceptors (rods and cones) and the other, newly discovered, based on photo-receptive ganglion cells which act as rudimentary visual brightness detectors.
Photochemistry
The functioning of a camera is often compared with the workings of the eye, mostly since both focus light from external objects in the field of view onto a light-sensitive medium. In the case of the camera, this medium is film or an electronic sensor; in the case of the eye, it is an array of visual receptors. With this simple geometrical similarity, based on the laws of optics, the eye functions as a transducer, as does a CCD camera.
In the visual system, retinal, technically called retinene1 or "retinaldehyde", is a light-sensitive molecule found in the rods and cones of the retina. Retinal is the fundamental structure involved in the transduction of light into visual signals, i.e. nerve impulses in the ocular system of the central nervous system. In the presence of light, the retinal molecule changes configuration and as a result, a nerve impulse is generated.
Optic nerve
The information about the image via the eye is transmitted to the brain along the optic nerve. Different populations of ganglion cells in the retina send information to the brain through the optic nerve. About 90% of the axons in the optic nerve go to the lateral geniculate nucleus in the thalamus. These axons originate from the M, P, and K ganglion cells in the retina, see above. This parallel processing is important for reconstructing the visual world; each type of information will go through a different route to perception. Another population sends information to the superior colliculus in the midbrain, which assists in controlling eye movements (saccades) as well as other motor responses.
A final population of photosensitive ganglion cells, containing melanopsin for photosensitivity, sends information via the retinohypothalamic tract to the pretectum (pupillary reflex), to several structures involved in the control of circadian rhythms and sleep such as the suprachiasmatic nucleus (the biological clock), and to the ventrolateral preoptic nucleus (a region involved in sleep regulation). A recently discovered role for photoreceptive ganglion cells is that they mediate conscious and unconscious vision – acting as rudimentary visual brightness detectors as shown in rodless coneless eyes.
Optic chiasm
The optic nerves from both eyes meet and cross at the optic chiasm, at the base of the hypothalamus of the brain. At this point, the information coming from both eyes is combined and then splits according to the visual field. The corresponding halves of the field of view (right and left) are sent to the left and right halves of the brain, respectively, to be processed. That is, the right side of primary visual cortex deals with the left half of the field of view from both eyes, and similarly for the left brain. A small region in the center of the field of view is processed redundantly by both halves of the brain.
Optic tract
Information from the right visual field (now on the left side of the brain) travels in the left optic tract. Information from the left visual field travels in the right optic tract. Each optic tract terminates in the lateral geniculate nucleus (LGN) in the thalamus.
Lateral geniculate nucleus
The lateral geniculate nucleus (LGN) is a sensory relay nucleus in the thalamus of the brain. The LGN consists of six layers in humans and other primates starting from catarrhines, including cercopithecidae and apes. Layers 1, 4, and 6 correspond to information from the contralateral (crossed) fibers of the nasal retina (temporal visual field); layers 2, 3, and 5 correspond to information from the ipsilateral (uncrossed) fibers of the temporal retina (nasal visual field).
Layer one contains M cells, which correspond to the M (magnocellular) cells of the optic nerve of the opposite eye and are concerned with depth or motion. Layers four and six of the LGN also connect to the opposite eye, but to the P cells (color and edges) of the optic nerve. By contrast, layers two, three and five of the LGN connect to the M cells and P (parvocellular) cells of the optic nerve for the same side of the brain as its respective LGN.
Spread out, the six layers of the LGN are the area of a credit card and about three times its thickness. The LGN is rolled up into two ellipsoids about the size and shape of two small birds' eggs. In between the six layers are smaller cells that receive information from the K cells (color) in the retina. The neurons of the LGN then relay the visual image to the primary visual cortex (V1) which is located at the back of the brain (posterior end) in the occipital lobe in and close to the calcarine sulcus. The LGN is not just a simple relay station, but it is also a center for processing; it receives reciprocal input from the cortical and subcortical layers and reciprocal innervation from the visual cortex.
Optic radiation
The optic radiations, one on each side of the brain, carry information from the thalamic lateral geniculate nucleus to layer 4 of the visual cortex. The P layer neurons of the LGN relay to V1 layer 4C β. The M layer neurons relay to V1 layer 4C α. The K layer neurons in the LGN relay to large neurons called blobs in layers 2 and 3 of V1.
There is a direct correspondence from an angular position in the visual field of the eye, all the way through the optic tract to a nerve position in V1 up to V4, i.e. the primary visual areas. After that, the visual pathway is roughly separated into a ventral and dorsal pathway.
Visual cortex
The visual cortex is responsible for processing the visual image. It lies at the rear of the brain (highlighted in the image), above the cerebellum. The region that receives information directly from the LGN is called the primary visual cortex (also called V1 and striate cortex). It creates a bottom-up saliency map of the visual field to guide attention or eye gaze to salient visual locations. Hence selection of visual input information by attention starts at V1 along the visual pathway.
Visual information then flows through a cortical hierarchy. These areas include V2, V3, V4 and area V5/MT. (The exact connectivity depends on the species of the animal.) These secondary visual areas (collectively termed the extrastriate visual cortex) process a wide variety of visual primitives. Neurons in V1 and V2 respond selectively to bars of specific orientations, or combinations of bars. These are believed to support edge and corner detection. Similarly, basic information about color and motion is processed here.
Heider, et al. (2002) found that neurons involving V1, V2, and V3 can detect stereoscopic illusory contours; they found that stereoscopic stimuli subtending up to 8° can activate these neurons.
Visual association cortex
As visual information passes forward through the visual hierarchy, the complexity of the neural representations increases. Whereas a V1 neuron may respond selectively to a line segment of a particular orientation in a particular retinotopic location, neurons in the lateral occipital complex respond selectively to a complete object (e.g., a figure drawing), and neurons in the visual association cortex may respond selectively to human faces, or to a particular object.
Along with this increasing complexity of neural representation may come a level of specialization of processing into two distinct pathways: the dorsal stream and the ventral stream (the Two Streams hypothesis, first proposed by Ungerleider and Mishkin in 1982). The dorsal stream, commonly referred to as the "where" stream, is involved in spatial attention (covert and overt), and communicates with regions that control eye movements and hand movements. More recently, this area has been called the "how" stream to emphasize its role in guiding behaviors to spatial locations. The ventral stream, commonly referred to as the "what" stream, is involved in the recognition, identification and categorization of visual stimuli.
However, there is still much debate about the degree of specialization within these two pathways, since they are in fact heavily interconnected.
Horace Barlow proposed the efficient coding hypothesis in 1961 as a theoretical model of sensory coding in the brain. Limitations in the applicability of this theory in the primary visual cortex (V1) motivated the V1 Saliency Hypothesis that V1 creates a bottom-up saliency map to guide attention exogenously. With attentional selection as a center stage, vision is seen as composed of encoding, selection, and decoding stages.
The default mode network is a network of brain regions that are active when an individual is awake and at rest. The visual system's default mode can be monitored during resting state fMRI:
Fox, et al. (2005) found that "the human brain is intrinsically organized into dynamic, anticorrelated functional networks", in which the visual system switches from resting state to attention.
In the parietal lobe, the lateral and ventral intraparietal cortex are involved in visual attention and saccadic eye movements. These regions are in the intraparietal sulcus (marked in red in the adjacent image).
Development
Infancy
Newborn infants have limited color perception. One study found that 74% of newborns can distinguish red, 36% green, 25% yellow, and 14% blue. After one month, performance "improved somewhat." Infant's eyes do not have the ability to accommodate. Pediatricians are able to perform non-verbal testing to assess visual acuity of a newborn, detect nearsightedness and astigmatism, and evaluate the eye teaming and alignment. Visual acuity improves from about 20/400 at birth to approximately 20/25 at 6 months of age. This happens because the nerve cells in the retina and brain that control vision are not fully developed.
Childhood and adolescence
Depth perception, focus, tracking and other aspects of vision continue to develop throughout early and middle childhood. From recent studies in the United States and Australia there is some evidence that the amount of time school aged children spend outdoors, in natural light, may have some impact on whether they develop myopia. The condition tends to get somewhat worse through childhood and adolescence, but stabilizes in adulthood. More prominent myopia (nearsightedness) and astigmatism are thought to be inherited. Children with this condition may need to wear glasses.
Adulthood
Vision is often one of the first senses affected by aging. A number of changes occur with aging:
Over time, the lens becomes yellowed and may eventually become brown, a condition known as brunescence or brunescent cataract. Although many factors contribute to yellowing, lifetime exposure to ultraviolet light and aging are two main causes.
The lens becomes less flexible, diminishing the ability to accommodate (presbyopia).
While a healthy adult pupil typically has a size range of 2–8 mm, with age the range gets smaller, trending towards a moderately small diameter.
On average tear production declines with age. However, there are a number of age-related conditions that can cause excessive tearing.
Other functions
Balance
Along with proprioception and vestibular function, the visual system plays an important role in the ability of an individual to control balance and maintain an upright posture. When these three conditions are isolated and balance is tested, it has been found that vision is the most significant contributor to balance, playing a bigger role than either of the two other intrinsic mechanisms. The clarity with which an individual can see his environment, as well as the size of the visual field, the susceptibility of the individual to light and glare, and poor depth perception play important roles in providing a feedback loop to the brain on the body's movement through the environment. Anything that affects any of these variables can have a negative effect on balance and maintaining posture. This effect has been seen in research involving elderly subjects when compared to young controls, in glaucoma patients compared to age matched controls, cataract patients pre and post surgery, and even something as simple as wearing safety goggles. Monocular vision (one eyed vision) has also been shown to negatively impact balance, which was seen in the previously referenced cataract and glaucoma studies, as well as in healthy children and adults.
According to Pollock et al. (2010) stroke is the main cause of specific visual impairment, most frequently visual field loss (homonymous hemianopia, a visual field defect). Nevertheless, evidence for the efficacy of cost-effective interventions aimed at these visual field defects is still inconsistent.
Clinical significance
Proper function of the visual system is required for sensing, processing, and understanding the surrounding environment. Difficulty in sensing, processing and understanding light input has the potential to adversely impact an individual's ability to communicate, learn and effectively complete routine tasks on a daily basis.
In children, early diagnosis and treatment of impaired visual system function is an important factor in ensuring that key social, academic and speech/language developmental milestones are met.
Cataract is clouding of the lens, which in turn affects vision. Although it may be accompanied by yellowing, clouding and yellowing can occur separately. This is typically a result of ageing, disease, or drug use.
Presbyopia is a visual condition that causes farsightedness. The eye's lens becomes too inflexible to accommodate to normal reading distance, focus tending to remain fixed at long distance.
Glaucoma is a type of blindness that begins at the edge of the visual field and progresses inward. It may result in tunnel vision. This typically involves the outer layers of the optic nerve, sometimes as a result of buildup of fluid and excessive pressure in the eye.
Scotoma is a type of blindness that produces a small blind spot in the visual field typically caused by injury in the primary visual cortex.
Homonymous hemianopia is a type of blindness that destroys one entire side of the visual field typically caused by injury in the primary visual cortex.
Quadrantanopia is a type of blindness that destroys only a part of the visual field typically caused by partial injury in the primary visual cortex. This is very similar to homonymous hemianopia, but to a lesser degree.
Prosopagnosia, or face blindness, is a brain disorder that produces an inability to recognize faces. This disorder often arises after damage to the fusiform face area.
Visual agnosia, or visual-form agnosia, is a brain disorder that produces an inability to recognize objects. This disorder often arises after damage to the ventral stream.
Other animals
Different species are able to see different parts of the light spectrum; for example, bees can see into the ultraviolet, while pit vipers can accurately target prey with their pit organs, which are sensitive to infrared radiation. The mantis shrimp possesses arguably the most complex visual system of any species. The eye of the mantis shrimp holds 16 color receptive cones, whereas humans only have three. The variety of cones enables them to perceive an enhanced array of colors as a mechanism for mate selection, avoidance of predators, and detection of prey. Swordfish also possess an impressive visual system. The eye of a swordfish can generate heat to better cope with detecting their prey at depths of 2000 feet. Certain one-celled microorganisms, the warnowiid dinoflagellates have eye-like ocelloids, with analogous structures for the lens and retina of the multi-cellular eye. The armored shell of the chiton Acanthopleura granulata is also covered with hundreds of aragonite crystalline eyes, named ocelli, which can form images.
Many fan worms, such as Acromegalomma interruptum which live in tubes on the sea floor of the Great Barrier Reef, have evolved compound eyes on their tentacles, which they use to detect encroaching movement. If movement is detected, the fan worms will rapidly withdraw their tentacles. Bok, et al., have discovered opsins and G proteins in the fan worm's eyes, which were previously only seen in simple ciliary photoreceptors in the brains of some invertebrates, as opposed to the rhabdomeric receptors in the eyes of most invertebrates.
Only higher primate Old World (African) monkeys and apes (macaques, apes, orangutans) have the same kind of three-cone photoreceptor color vision humans have, while lower primate New World (South American) monkeys (spider monkeys, squirrel monkeys, cebus monkeys) have a two-cone photoreceptor kind of color vision.
Biologists have determined that humans have extremely good vision compared to the overwhelming majority of animals, particularly in daylight, surpassed only by a few large species of predatory birds. Other animals such as dogs are thought to rely more on senses other than vision, which in turn may be better developed than in humans.
History
In the second half of the 19th century, many motifs of the nervous system were identified such as the neuron doctrine and brain localization, which related to the neuron being the basic unit of the nervous system and functional localisation in the brain, respectively. These would become tenets of the fledgling neuroscience and would support further understanding of the visual system.
The notion that the cerebral cortex is divided into functionally distinct cortices now known to be responsible for capacities such as touch (somatosensory cortex), movement (motor cortex), and vision (visual cortex), was first proposed by Franz Joseph Gall in 1810. Evidence for functionally distinct areas of the brain (and, specifically, of the cerebral cortex) mounted throughout the 19th century with discoveries by Paul Broca of the language center (1861), and Gustav Fritsch and Eduard Hitzig of the motor cortex (1871). Based on selective damage to parts of the brain and the functional effects of the resulting lesions, David Ferrier proposed that visual function was localized to the parietal lobe of the brain in 1876. In 1881, Hermann Munk more accurately located vision in the occipital lobe, where the primary visual cortex is now known to be.
In 2014, a textbook "Understanding vision: theory, models, and data" illustrates how to link neurobiological data and visual behavior/psychological data through theoretical principles and computational models.
| Biology and health sciences | Nervous system | null |
305225 | https://en.wikipedia.org/wiki/Manganese%20dioxide | Manganese dioxide | Manganese dioxide is the inorganic compound with the formula . This blackish or brown solid occurs naturally as the mineral pyrolusite, which is the main ore of manganese and a component of manganese nodules. The principal use for is for dry-cell batteries, such as the alkaline battery and the zinc–carbon battery. is also used as a pigment and as a precursor to other manganese compounds, such as . It is used as a reagent in organic synthesis, for example, for the oxidation of allylic alcohols. has an α-polymorph that can incorporate a variety of atoms (as well as water molecules) in the "tunnels" or "channels" between the manganese oxide octahedra. There is considerable interest in as a possible cathode for lithium-ion batteries.
Structure
Several polymorphs of are claimed, as well as a hydrated form. Like many other dioxides, crystallizes in the rutile crystal structure (this polymorph is called pyrolusite or ), with three-coordinate oxide anions and octahedral metal centres. is characteristically nonstoichiometric, being deficient in oxygen. The complicated solid-state chemistry of this material is relevant to the lore of "freshly prepared" in organic synthesis. The α-polymorph of has a very open structure with "channels", which can accommodate metal ions such as silver or barium. is often called hollandite, after a closely related mineral.
Production
Naturally occurring manganese dioxide contains impurities and a considerable amount of manganese(III) oxide. Production of batteries and ferrite (two of the primary uses of manganese dioxide) requires high purity manganese dioxide. Batteries require "electrolytic manganese dioxide" while ferrites require "chemical manganese dioxide".
Chemical manganese dioxide
One method starts with natural manganese dioxide and converts it using dinitrogen tetroxide and water to a manganese(II) nitrate solution. Evaporation of the water leaves the crystalline nitrate salt. At temperatures of 400 °C, the salt decomposes, releasing and leaving a residue of purified manganese dioxide. These two steps can be summarized as:
+
In another process, manganese dioxide is carbothermically reduced to manganese(II) oxide which is dissolved in sulfuric acid. The filtered solution is treated with ammonium carbonate to precipitate . The carbonate is calcined in air to give a mixture of manganese(II) and manganese(IV) oxides. To complete the process, a suspension of this material in sulfuric acid is treated with sodium chlorate. Chloric acid, which forms in situ, converts any Mn(III) and Mn(II) oxides to the dioxide, releasing chlorine as a by-product.
Lastly, the action of potassium permanganate over manganese sulfate crystals produces the desired oxide.
2 + 3 + 2 → 5 + + 2
Electrolytic manganese dioxide
Electrolytic manganese dioxide (EMD) is used in zinc–carbon batteries together with zinc chloride and ammonium chloride. EMD is commonly used in zinc manganese dioxide rechargeable alkaline (Zn RAM) cells also. For these applications, purity is extremely important. EMD is produced in a similar fashion as electrolytic tough pitch (ETP) copper: The manganese dioxide is dissolved in sulfuric acid (sometimes mixed with manganese sulfate) and subjected to a current between two electrodes. The MnO2 dissolves, enters solution as the sulfate, and is deposited on the anode.
Reactions
The important reactions of are associated with its redox, both oxidation and reduction.
Reduction
is the principal precursor to ferromanganese and related alloys, which are widely used in the steel industry. The conversions involve carbothermal reduction using coke:
+ 2 C → Mn + 2 CO
The key redox reactions of in batteries is the one-electron reduction:
+ e− + → MnO(OH)
catalyses several reactions that form . In a classical laboratory demonstration, heating a mixture of potassium chlorate and manganese dioxide produces oxygen gas. Manganese dioxide also catalyses the decomposition of hydrogen peroxide to oxygen and water:
2 → 2 +
Manganese dioxide decomposes above about 530 °C to manganese(III) oxide and oxygen. At temperatures close to 1000 °C, the mixed-valence compound forms. Higher temperatures give MnO, which is reduced only with difficulty.
Hot concentrated sulfuric acid reduces to manganese(II) sulfate:
2 + 2 → 2 + + 2
The reaction of hydrogen chloride with was used by Carl Wilhelm Scheele in the original isolation of chlorine gas in 1774:
+ 4 HCl → + + 2
As a source of hydrogen chloride, Scheele treated sodium chloride with concentrated sulfuric acid.
Eo ((s) + 4 + 2 e− Mn2+ + 2 ) = +1.23 V
Eo ((g) + 2 e− 2 Cl−) = +1.36 V
The reaction would not be expected to proceed, based on the standard electrode potentials, but is favoured by the extremely high acidity and the evolution (and removal) of gaseous chlorine.
This reaction is also a convenient way to remove the manganese dioxide precipitate from the ground glass joints after running a reaction (for example, an oxidation with potassium permanganate).
Oxidation
Heating a mixture of KOH and in air gives green potassium manganate:
2 + 4 KOH + → 2 + 2
Potassium manganate is the precursor to potassium permanganate, a common oxidant.
Occurrence and applications
The predominant application of is as a component of dry cell batteries: alkaline batteries and so called Leclanché cell, or zinc–carbon batteries. Approximately 500,000 tonnes are consumed for this application annually. Other industrial applications include the use of as an inorganic pigment in ceramics and in glassmaking. It is also used in water treatment applications.
Prehistory
Excavations at the Pech-de-l'Azé cave site in southwestern France have yielded blocks of manganese dioxide writing tools, which date back 50,000 years and have been attributed to Neanderthals . Scientists have conjectured that Neanderthals used this mineral for body decoration, but there are many other readily available minerals that are more suitable for that purpose. Heyes et al. (in 2016) determined that the manganese dioxide lowers the combustion temperatures for wood from above 350°C (662°F) to 250°C (482°F), making fire making much easier and this is likely to be the purpose of the blocks.
Organic synthesis
A specialized use of manganese dioxide is as oxidant in organic synthesis. The effectiveness of the reagent depends on the method of preparation, a problem that is typical for other heterogeneous reagents where surface area, among other variables, is a significant factor. The mineral pyrolusite makes a poor reagent. Usually, however, the reagent is generated in situ by treatment of an aqueous solution with a Mn(II) salt, typically the sulfate. oxidizes allylic alcohols to the corresponding aldehydes or ketones:
cis-RCH= + → cis-RCH=CHCHO + MnO +
The configuration of the double bond is conserved in the reaction. The corresponding acetylenic alcohols are also suitable substrates, although the resulting propargylic aldehydes can be quite reactive. Benzylic and even unactivated alcohols are also good substrates. 1,2-Diols are cleaved by to dialdehydes or diketones. Otherwise, the applications of are numerous, being applicable to many kinds of reactions including amine oxidation, aromatization, oxidative coupling, and thiol oxidation.
Microbiology
In Geobacteraceae sp., MnO2 functions as an electron acceptor coupled to the oxidation of organic compounds. This theme has implications for bioremediation.
| Physical sciences | Oxide salts | Chemistry |
305286 | https://en.wikipedia.org/wiki/Legume | Legume | Legumes () are plants in the family Fabaceae (or Leguminosae), or the fruit or seeds of such plants. When used as a dry grain for human consumption, the seeds are also called pulses. Legumes are grown agriculturally, primarily for human consumption, but also as livestock forage and silage, and as soil-enhancing green manure. Well-known legumes include beans, chickpeas, peanuts, lentils, lupins, mesquite, carob, tamarind, alfalfa, and clover. Legumes produce a botanically unique type of fruit – a simple dry fruit that develops from a simple carpel and usually dehisces (opens along a seam) on two sides.
Most legumes have symbiotic nitrogen-fixing bacteria, Rhizobia, in structures called root nodules. Some of the fixed nitrogen becomes available to later crops, so legumes play a key role in crop rotation.
Terminology
The term pulse, as used by the United Nations' Food and Agriculture Organization (FAO), is reserved for legume crops harvested solely for the dry seed. This excludes green beans and green peas, which are considered vegetable crops. Also excluded are seeds that are mainly grown for oil extraction (oilseeds like soybeans and peanuts), and seeds which are used exclusively for sowing forage (clovers, alfalfa). However, in common usage, these distinctions are not always clearly made, and many of the varieties used for dried pulses are also used for green vegetables, with their beans in pods while young.
Some Fabaceae, such as Scotch broom and other Genisteae, are leguminous but are usually not called legumes by farmers, who tend to restrict that term to food crops.
The FAO recognizes 11 primary pulses, excluding green vegetable legumes (e.g. green peas) and legumes used mainly for oil extraction (e.g., soybeans and groundnuts) or used only as seed (e.g., clover and alfalfa).
Dry beans (FAOSTAT code 0176, Phaseolus spp. including several species now in Vigna)
Kidney bean, navy bean, pinto bean, black turtle bean, haricot bean (Phaseolus vulgaris)
Lima bean, butter bean (Phaseolus lunatus)
Adzuki bean, azuki bean (Vigna angularis)
Mung bean, golden gram, green gram (Vigna radiata)
Black gram, urad (Vigna mungo)
Scarlet runner bean (Phaseolus coccineus)
Ricebean (Vigna umbellata)
Moth bean (Vigna aconitifolia)
Tepary bean (Phaseolus acutifolius)
Dry broad beans (code 0181, Vicia faba)
Horse bean (Vicia faba equina)
Broad bean (Vicia faba)
Field bean (Vicia faba)
Dry peas (code 0187, Pisum spp.)
Garden pea (Pisum sativum var. sativum)
Protein pea (Pisum sativum var. arvense)
Chickpea, garbanzo, Bengal gram (code 0191, Cicer arietinum)
Dry cowpea, black-eyed pea, blackeye bean (code 0195, Vigna unguiculata)
Pigeon pea, Arhar/Toor, cajan pea, Congo bean, gandules (code 0197, Cajanus cajan)
Lentil (code 0201, Lens culinaris)
Bambara groundnut, earth pea (code 0203, Vigna subterranea)
Vetch, common vetch (code 0205, Vicia sativa)
Lupins (code 0210, Lupinus spp.)
Pulses NES (code 0211), Minor pulses, including:
Lablab, hyacinth bean (Lablab purpureus)
Jack bean (Canavalia ensiformis), sword bean (Canavalia gladiata)
Winged bean (Psophocarpus tetragonolobus)
Velvet bean, cowitch (Mucuna pruriens var. utilis)
Yam bean (Pachyrhizus erosus)
Distribution
Legumes are widely distributed as the third-largest land plant family in terms of number of species, behind only the Orchidaceae and Asteraceae, with about 751 genera and some 19,000 known species, constituting about seven percent of flowering plant species.
Ecology
Nitrogen fixation
Many legumes contain symbiotic bacteria called Rhizobia within root nodules of their root systems (plants belonging to the genus Styphnolobium are one exception to this rule). These bacteria have the special ability of fixing nitrogen from atmospheric, molecular nitrogen (N2) into ammonia (NH3). The chemical reaction is:
Ammonia is converted to another form, ammonium (), usable by (some) plants, by the following reaction:
This arrangement means that the root nodules are sources of nitrogen for legumes, making them relatively rich in plant proteins. All proteins contain nitrogenous amino acids. Nitrogen is therefore a necessary ingredient in the production of proteins. Hence, legumes are among the best sources of plant protein.
When a legume plant dies in the field, for example following the harvest, all of its remaining nitrogen, incorporated into amino acids inside the remaining plant parts, is released back into the soil. In the soil, the amino acids are converted to nitrate (), making the nitrogen available to other plants, thereby serving as fertilizer for future crops.
In many traditional and organic farming practices, crop rotation or polyculture involving legumes is common. By alternating between legumes and non-legumes, or by growing both together for part of the growing season, the field can receive a sufficient amount of nitrogenous compounds to produce a good result without adding nitrogenous fertilizer. Legumes are often used as green manure.
Sri Lanka developed the polyculture practice known as coconut-soybean intercropping. Grain legumes are grown in coconut (Cocos nuficera) groves in two ways: intercropping or as a cash crop. These are grown mainly for their protein, vegetable oil and ability to uphold soil fertility. However, continuous cropping after 3–4 years decrease grain yields significantly.
Pests and diseases
A common pest of grain legumes that is noticed in the tropical and subtropical Asia, Africa, Australia and Oceania are minuscule flies that belong to the family Agromyzidae, dubbed "bean flies". They are considered to be the most destructive. The host range of these flies is very wide amongst cultivated legumes. Infestation of plants starts from germination through to harvest, and they can destroy an entire crop in early stage. Black bean aphids are a serious pest to broad beans and other beans. Common hosts for this pest are fathen, thistle and dock. Pea weevil and bean weevil damage leaf margins leaving characteristics semi-circular notches. Stem nematodes are very widespread but will be found more frequently in areas where host plants are grown.
Common legume diseases include anthracnose, caused by Colletotrichum trifolii; common leaf spot caused by Pseudomonas syringae pv. syringae; crown wart caused by Physoderma alfalfae; downy mildew caused by Peronospora trifoliorum; fusarium root rot caused by Fusarium spp.; rust caused by Uromyces striatus; sclerotina crown and stem rot caused by Sclerotinia trifoliorum; Southern blight caused by Sclerotium rolfsii; pythium (browning) root rot caused by Pythium spp.; fusarium wilt caused by Fusarium oxysporum; root knot caused by Meloidogyne hapla. These are all classified as biotic problems.
Abiotic problems include nutrient deficiencies, (nitrogen, phosphorus, potassium, copper, magnesium, manganese, boron, zinc), pollutants (air, water, soil, pesticide injury, fertilizer burn), toxic concentration of minerals, and unfavorable growth conditions.
Storage
Seed viability decreases with longer storage time. Studies of vetch, broad beans, and peas show that they last about 5 years in storage. Environmental factors that are important in influencing germination are relative humidity and temperature. Two rules apply to moisture content between 5 and 14 percent: the life of the seed will last longer if the storage temperature is reduced by 5 degree Celsius. Secondly, the storage moisture content will decrease if temperature is reduced by 1 degree Celsius.
Uses
Cultivated legumes encompass a diverse range of agricultural classifications, spanning forage, grain, flowering, pharmaceutical/industrial, fallow/green manure, and timber categories. A notable characteristic of many commercially cultivated legume species is their versatility, often assuming multiple roles concurrently. The extent of these roles is contingent upon the stage of maturity at which they are harvested.
Human consumption
Grain legumes are cultivated for their seeds, for humans and animals to eat, or for oils for industrial uses. Grain legumes include beans, lentils, lupins, peas, and peanuts.
Legumes are a key ingredient in vegan meat and dairy substitutes. They are growing in use as a plant-based protein source in the world marketplace. Products containing legumes grew by 39% in Europe between 2013 and 2017.
There is a common misconception that adding salt before cooking prevents them from cooking through. Legumes may not soften because they are old, or because of hard water or acidic ingredients in the pot; salting before cooking results in better seasoning.
Nutritional value
Legumes are a significant source of protein, dietary fibre, carbohydrates, and dietary minerals; for example, a 100 gram serving of cooked chickpeas contains 18 percent of the Daily Value (DV) for protein, 30 percent DV for dietary fiber, 43 percent DV for folate and 52 percent DV for manganese.
Legumes are an excellent source of resistant starch; this is broken down by bacteria in the large intestine to produce short-chain fatty acids (such as butyrate) used by intestinal cells for food energy.
Forage
Forage legumes are of two broad types. Some, like alfalfa, clover, vetch (Vicia), stylo (Stylosanthes), or Arachis, are sown in pasture and grazed by livestock. Others, such as Leucaena or Albizia, are woody shrubs or trees that are either broken down by livestock or regularly cut by humans to provide fodder. Legume-based feeds improve animal performance over a diet of perennial grasses. Factors include larger consumption, faster digestion, and higher feed conversion rate.
The type of crop grown for animal rearing depends on the farming system. In cattle rearing, legume trees such as Gliricidia sepium can be planted along edges of fields to provide shade for cattle, the leaves and bark are often eaten by cattle. Green manure can be grown between harvesting the main crop and the planting of the next crop.
Other uses
Legume species grown for their flowers include lupins, which are farmed commercially for their blooms as well as being popular in gardens worldwide. Industrially farmed legumes include Indigofera and Acacia species, which are cultivated for dye and natural gum production, respectively. Fallow or green manure legume species are cultivated to be tilled back into the soil in order to exploit the high levels of captured atmospheric nitrogen found in the roots of most legumes. Numerous legumes farmed for this purpose include Leucaena, Cyamopsis, and Sesbania species. Various legume species are farmed for timber production worldwide, including numerous Acacia species and Castanospermum australe.
Some legume trees, like the honey locust (Gleditsia) can be used in agroforestry. Others, including the black locust (Robinia pseudoacacia), Kentucky coffeetree (Gymnocladus dioicus), Laburnum, and the woody climbing vine Wisteria, have poisonous elements.
History
Neanderthals and early modern humans used wild pulses when cooking meals 70,000 to 40,000 years ago. Traces of pulse production have been found around the Ravi River (Punjab), the seat of the Indus Valley civilisation, from 3300 BC. Meanwhile, evidence of lentil cultivation has also been found in Egyptian pyramids and cuneiform recipes. Dry pea seeds have been discovered in a Swiss village that are believed to date back to the Stone Age. Archaeological evidence suggests that these peas must have been grown in the eastern Mediterranean and Mesopotamian regions at least 5,000 years ago and in Britain as early as the 11th century. The soybean was domesticated around 5,000 years ago in China from a descendant of the wild vine Glycine soja.
The oldest-known domesticated beans in the Americas were found in Guitarrero Cave, an archaeological site in Peru, and dated to around the second millennium BCE. Genetic analyses of the common bean Phaseolus show that it originated in Mesoamerica, and subsequently spread southward, along with maize and squash, traditional companion crops. In the United States, the domesticated soybean was introduced in 1770 by Benjamin Franklin after he sent seeds to Philadelphia from France.
International Year of Pulses
The International Year of Pulses 2016 was declared by the Sixty-eighth session of the United Nations General Assembly. The Food and Agriculture Organization of the United Nations was nominated to facilitate the implementation of the year in collaboration with governments, relevant organizations, non-governmental organizations and other relevant stakeholders. Its aim was to heighten public awareness of the nutritional benefits of pulses as part of sustainable food production aimed towards food security and nutrition. The year created an opportunity to encourage connections throughout the food chain that would better use pulse-based proteins, further global production of pulses, better use crop rotations and address challenges in the global trade of pulses.
| Biology and health sciences | Pulses | Plants |
305303 | https://en.wikipedia.org/wiki/Multiset | Multiset | In mathematics, a multiset (or bag, or mset) is a modification of the concept of a set that, unlike a set, allows for multiple instances for each of its elements. The number of instances given for each element is called the multiplicity of that element in the multiset. As a consequence, an infinite number of multisets exist that contain only elements and , but vary in the multiplicities of their elements:
The set contains only elements and , each having multiplicity 1 when is seen as a multiset.
In the multiset , the element has multiplicity 2, and has multiplicity 1.
In the multiset , and both have multiplicity 3.
These objects are all different when viewed as multisets, although they are the same set, since they all consist of the same elements. As with sets, and in contrast to tuples, the order in which elements are listed does not matter in discriminating multisets, so and denote the same multiset. To distinguish between sets and multisets, a notation that incorporates square brackets is sometimes used: the multiset can be denoted by .
The cardinality or "size" of a multiset is the sum of the multiplicities of all its elements. For example, in the multiset the multiplicities of the members , , and are respectively 2, 3, and 1, and therefore the cardinality of this multiset is 6.
Nicolaas Govert de Bruijn coined the word multiset in the 1970s, according to Donald Knuth. However, the concept of multisets predates the coinage of the word multiset by many centuries. Knuth himself attributes the first study of multisets to the Indian mathematician Bhāskarāchārya, who described permutations of multisets around 1150. Other names have been proposed or used for this concept, including list, bunch, bag, heap, sample, weighted set, collection, and suite.
History
Wayne Blizard traced multisets back to the very origin of numbers, arguing that "in ancient times, the number n was often represented by a collection of n strokes, tally marks, or units." These and similar collections of objects can be regarded as multisets, because strokes, tally marks, or units are considered indistinguishable. This shows that people implicitly used multisets even before mathematics emerged.
Practical needs for this structure have caused multisets to be rediscovered several times, appearing in literature under different names. For instance, they were important in early AI languages, such as QA4, where they were referred to as bags, a term attributed to Peter Deutsch. A multiset has been also called an aggregate, heap, bunch, sample, weighted set, occurrence set, and fireset (finitely repeated element set).
Although multisets were used implicitly from ancient times, their explicit exploration happened much later. The first known study of multisets is attributed to the Indian mathematician Bhāskarāchārya circa 1150, who described permutations of multisets. The work of Marius Nizolius (1498–1576) contains another early reference to the concept of multisets. Athanasius Kircher found the number of multiset permutations when one element can be repeated. Jean Prestet published a general rule for multiset permutations in 1675. John Wallis explained this rule in more detail in 1685.
Multisets appeared explicitly in the work of Richard Dedekind.
Other mathematicians formalized multisets and began to study them as precise mathematical structures in the 20th century. For example, Hassler Whitney (1933) described generalized sets ("sets" whose characteristic functions may take any integer value: positive, negative or zero). Monro (1987) investigated the category Mul of multisets and their morphisms, defining a multiset as a set with an equivalence relation between elements "of the same sort", and a morphism between multisets as a function that respects sorts. He also introduced a multinumber : a function f (x) from a multiset to the natural numbers, giving the multiplicity of element x in the multiset. Monro argued that the concepts of multiset and multinumber are often mixed indiscriminately, though both are useful.
Examples
One of the simplest and most natural examples is the multiset of prime factors of a natural number . Here the underlying set of elements is the set of prime factors of . For example, the number 120 has the prime factorization
which gives the multiset .
A related example is the multiset of solutions of an algebraic equation. A quadratic equation, for example, has two solutions. However, in some cases they are both the same number. Thus the multiset of solutions of the equation could be , or it could be . In the latter case it has a solution of multiplicity 2. More generally, the fundamental theorem of algebra asserts that the complex solutions of a polynomial equation of degree always form a multiset of cardinality .
A special case of the above are the eigenvalues of a matrix, whose multiplicity is usually defined as their multiplicity as roots of the characteristic polynomial. However two other multiplicities are naturally defined for eigenvalues, their multiplicities as roots of the minimal polynomial, and the geometric multiplicity, which is defined as the dimension of the kernel of (where is an eigenvalue of the matrix ). These three multiplicities define three multisets of eigenvalues, which may be all different: Let be a matrix in Jordan normal form that has a single eigenvalue. Its multiplicity is , its multiplicity as a root of the minimal polynomial is the size of the largest Jordan block, and its geometric multiplicity is the number of Jordan blocks.
Definition
A multiset may be formally defined as an ordered pair where is the underlying set of the multiset, formed from its distinct elements, and is a function from to the set of positive integers, giving the multiplicity – that is, the number of occurrences – of the element in the multiset as the number .
(It is also possible to allow multiplicity 0 or , especially when considering submultisets. This article is restricted to finite, positive multiplicities.)
Representing the function by its graph (the set of ordered pairs ) allows for writing the multiset as }, and the multiset as }. This notation is however not commonly used; more compact notations are employed.
If is a finite set, the multiset is often represented as
sometimes simplified to
where upper indices equal to 1 are omitted. For example, the multiset may be written or If the elements of the multiset are numbers, a confusion is possible with ordinary arithmetic operations; those normally can be excluded from the context. On the other hand, the latter notation is coherent with the fact that the prime factorization of a positive integer is a uniquely defined multiset, as asserted by the fundamental theorem of arithmetic. Also, a monomial is a multiset of indeterminates; for example, the monomial x3y2 corresponds to the multiset .
A multiset corresponds to an ordinary set if the multiplicity of every element is 1. An indexed family , where varies over some index set I, may define a multiset, sometimes written . In this view the underlying set of the multiset is given by the image of the family, and the multiplicity of any element is the number of index values such that . In this article the multiplicities are considered to be finite, so that no element occurs infinitely many times in the family; even in an infinite multiset, the multiplicities are finite numbers.
It is possible to extend the definition of a multiset by allowing multiplicities of individual elements to be infinite cardinals instead of positive integers, but not all properties carry over to this generalization.
Basic properties and operations
Elements of a multiset are generally taken in a fixed set , sometimes called a universe, which is often the set of natural numbers. An element of that does not belong to a given multiset is said to have a multiplicity 0 in this multiset. This extends the multiplicity function of the multiset to a function from to the set of non-negative integers. This defines a one-to-one correspondence between these functions and the multisets that have their elements in .
This extended multiplicity function is commonly called simply the multiplicity function, and suffices for defining multisets when the universe containing the elements has been fixed. This multiplicity function is a generalization of the indicator function of a subset, and shares some properties with it.
The support of a multiset in a universe is the underlying set of the multiset. Using the multiplicity function , it is characterized as
A multiset is finite if its support is finite, or, equivalently, if its cardinality
is finite. The empty multiset is the unique multiset with an empty support (underlying set), and thus a cardinality 0.
The usual operations of sets may be extended to multisets by using the multiplicity function, in a similar way to using the indicator function for subsets. In the following, and are multisets in a given universe , with multiplicity functions and
Inclusion: is included in , denoted , if
Union: the union (called, in some contexts, the maximum or lowest common multiple) of and is the multiset with multiplicity function
Intersection: the intersection (called, in some contexts, the infimum or greatest common divisor) of and is the multiset with multiplicity function
Sum: the sum of and is the multiset with multiplicity function It may be viewed as a generalization of the disjoint union of sets. It defines a commutative monoid structure on the finite multisets in a given universe. This monoid is a free commutative monoid, with the universe as a basis.
Difference: the difference of and is the multiset with multiplicity function
Two multisets are disjoint if their supports are disjoint sets. This is equivalent to saying that their intersection is the empty multiset or that their sum equals their union.
There is an inclusion–exclusion principle for finite multisets (similar to the one for sets), stating that a finite union of finite multisets is the difference of two sums of multisets: in the first sum we consider all possible intersections of an odd number of the given multisets, while in the second sum we consider all possible intersections of an even number of the given multisets.
Counting multisets
The number of multisets of cardinality , with elements taken from a finite set of cardinality , is sometimes called the multiset coefficient or multiset number. This number is written by some authors as , a notation that is meant to resemble that of binomial coefficients; it is used for instance in (Stanley, 1997), and could be pronounced " multichoose " to resemble " choose " for Like the binomial distribution that involves binomial coefficients, there is a negative binomial distribution in which the multiset coefficients occur. Multiset coefficients should not be confused with the unrelated multinomial coefficients that occur in the multinomial theorem.
The value of multiset coefficients can be given explicitly as
where the second expression is as a binomial coefficient; many authors in fact avoid separate notation and just write binomial coefficients. So, the number of such multisets is the same as the number of subsets of cardinality of a set of cardinality . The analogy with binomial coefficients can be stressed by writing the numerator in the above expression as a rising factorial power
to match the expression of binomial coefficients using a falling factorial power:
For example, there are 4 multisets of cardinality 3 with elements taken from the set of cardinality 2 (, ), namely , , , . There are also 4 subsets of cardinality 3 in the set of cardinality 4 (), namely , , , .
One simple way to prove the equality of multiset coefficients and binomial coefficients given above involves representing multisets in the following way. First, consider the notation for multisets that would represent (6 s, 2 s, 3 s, 7 s) in this form:
This is a multiset of cardinality made of elements of a set of cardinality . The number of characters including both dots and vertical lines used in this notation is . The number of vertical lines is 4 − 1. The number of multisets of cardinality 18 is then the number of ways to arrange the vertical lines among the 18 + 4 − 1 characters, and is thus the number of subsets of cardinality 4 − 1 of a set of cardinality . Equivalently, it is the number of ways to arrange the 18 dots among the characters, which is the number of subsets of cardinality 18 of a set of cardinality . This is
thus is the value of the multiset coefficient and its equivalencies:
From the relation between binomial coefficients and multiset coefficients, it follows that the number of multisets of cardinality in a set of cardinality can be written
Additionally,
Recurrence relation
A recurrence relation for multiset coefficients may be given as
with
The above recurrence may be interpreted as follows.
Let be the source set. There is always exactly one (empty) multiset of size 0, and if there are no larger multisets, which gives the initial conditions.
Now, consider the case in which . A multiset of cardinality with elements from might or might not contain any instance of the final element . If it does appear, then by removing once, one is left with a multiset of cardinality of elements from , and every such multiset can arise, which gives a total of
possibilities.
If does not appear, then our original multiset is equal to a multiset of cardinality with elements from , of which there are
Thus,
Generating series
The generating function of the multiset coefficients is very simple, being
As multisets are in one-to-one correspondence with monomials, is also the number of monomials of degree in indeterminates. Thus, the above series is also the Hilbert series of the polynomial ring
As is a polynomial in , it and the generating function are well defined for any complex value of .
Generalization and connection to the negative binomial series
The multiplicative formula allows the definition of multiset coefficients to be extended by replacing by an arbitrary number (negative, real, or complex):
With this definition one has a generalization of the negative binomial formula (with one of the variables set to 1), which justifies calling the negative binomial coefficients:
This Taylor series formula is valid for all complex numbers α and X with . It can also be interpreted as an identity of formal power series in X, where it actually can serve as definition of arbitrary powers of series with constant coefficient equal to 1; the point is that with this definition all identities hold that one expects for exponentiation, notably
and formulas such as these can be used to prove identities for the multiset coefficients.
If is a nonpositive integer , then all terms with are zero, and the infinite series becomes a finite sum. However, for other values of , including positive integers and rational numbers, the series is infinite.
Applications
Multisets have various applications. They are becoming fundamental in combinatorics. Multisets have become an important tool in the theory of relational databases, which often uses the synonym bag. For instance, multisets are often used to implement relations in database systems. In particular, a table (without a primary key) works as a multiset, because it can have multiple identical records. Similarly, SQL operates on multisets and returns identical records. For instance, consider "SELECT name from Student". In the case that there are multiple records with name "Sara" in the student table, all of them are shown. That means the result of an SQL query is a multiset; if the result were instead a set, the repetitive records in the result set would have been eliminated. Another application of multisets is in modeling multigraphs. In multigraphs there can be multiple edges between any two given vertices. As such, the entity that specifies the edges is a multiset, and not a set.
There are also other applications. For instance, Richard Rado used multisets as a device to investigate the properties of families of sets. He wrote, "The notion of a set takes no account of multiple occurrence of any one of its members, and yet it is just this kind of information that is frequently of importance. We need only think of the set of roots of a polynomial f (x) or the spectrum of a linear operator."
Generalizations
Different generalizations of multisets have been introduced, studied and applied to solving problems.
Real-valued multisets (in which multiplicity of an element can be any real number)
Fuzzy multisets
Rough multisets
Hybrid sets
Multisets whose multiplicity is any real-valued step function
Soft multisets
Soft fuzzy multisets
Named sets (unification of all generalizations of sets)
| Mathematics | Set theory | null |
305313 | https://en.wikipedia.org/wiki/Viscacha | Viscacha | Viscacha or vizcacha (, ) are rodents of two genera (Lagidium and Lagostomus) in the family Chinchillidae. They are native to South America and convergently resemble rabbits.
The five extant species of viscacha are:
The plains viscacha (Lagostomus maximus), a resident of the Pampas of Argentina, is easily differentiated from other viscachas by black and gray mustache-like facial markings. This species lives colonially in warrens of 10 to over 100. It is very vocal and emits alarm calls. The plains viscacha can strip grassland used to graze livestock; this has caused ranchers to consider the rodent a pest species.
Lagidium ahuacaense is a newly described species of mountain viscacha from the Ecuadorian Andes.
The northern viscacha (Lagidium peruanum) is native to the Peruvian Andes at elevations between the tree line and the snow line. It is dorsally gray or brown in color, with a bushy tail and long, furry ears. This species lives in large colonies separated into individual family units, like an apartment complex. It eats a wide range of plant matter, settling for almost anything it can find growing in the harsh, rocky environment.
The southern viscacha (Lagidium viscacia), also called mountain viscacha, is similar to the northern viscacha, but its pelage is more red in color. It lives in similar habitat in the Andes.
Wolffsohn's viscacha (Lagidium wolffsohni) is rarer than the other four species.
| Biology and health sciences | Rodents | Animals |
305399 | https://en.wikipedia.org/wiki/Porphyry%20%28geology%29 | Porphyry (geology) | Porphyry ( ) is any of various granites or igneous rocks with coarse-grained crystals such as feldspar or quartz dispersed in a fine-grained silicate-rich, generally aphanitic matrix or groundmass. In its non-geologic, traditional use, the term porphyry usually refers to the purple-red form of this stone, valued for its appearance, but other colours of decorative porphyry are also used such as "green", "black" and "grey".
The term porphyry is from the Ancient Greek (), meaning "purple". Purple was the colour of royalty, and the Roman "imperial porphyry" was a deep purple igneous rock with large crystals of plagioclase. Some authors claimed the rock was the hardest known in antiquity. Thus porphyry was prized for monuments and building projects in Imperial Rome and thereafter.
Subsequently, the name was given to any igneous rocks with large crystals. The adjective porphyritic now refers to a certain texture of igneous rock regardless of its chemical and mineralogical composition or its color. Its chief characteristic is a large difference in size between the tiny matrix crystals and the much larger phenocrysts. Porphyries may be aphanites or phanerites, that is, the groundmass may have microscopic crystals as in basalt, or crystals easily distinguishable with the eye, as in granite.
Formation
Most igneous rocks have some degree of porphyritic texture. This is because most magma from which igneous rock solidifies is produced by partial melting of a mixture of different minerals. At first the mixed melt slowly cools deep in the crust. The magma begins crystallizing, the highest melting point minerals closest to the overall composition first, in a process called fractional crystallization. This forms phenocrysts, which usually have plenty of room for growth, and form large, well-shaped crystals with characteristic crystal faces (euhedral crystals). If they are different in density to the remaining melt, these phenocrysts usually settle out of solution, eventually creating cumulates; however if the partially crystallized magma is then erupted to the surface as a lava, the remainder of the melt is quickly cooled around the phenocrysts and crystallizes much more rapidly to form a very fine-grained or glassy matrix.
Porphyry can also form even from magma that completely solidifies while still underground. The groundmass will be visibly crystalline, though not as large as the phenocrysts. The crystallization of the phenocrysts during fractional crystallization changes the composition of the remaining liquid magma, moving it closer to the eutectic point, with a mixed composition of minerals. As the temperature continues to decrease, this point is reached, and the rock is entirely solidified. The simultaneous crystallization of the remaining minerals produces the finer-grained matrix surrounding the phenocrysts, as they crowd each other out.
The significance of porphyritic texture as an indication that magma forms through different stages of cooling was first recognized by the Canadian geologist, Norman L. Bowen, in 1928.
Porphyritic texture is particularly common in andesite, with the most prominent phenocrysts typically composed of plagioclase feldspar. Plagioclase has almost the same density as basaltic magma, so plagioclase phenocrysts are likely to remain suspended in the magma rather than settling out.
Rhomb porphyry
Rhomb porphyry is a volcanic rock with gray-white large porphyritic rhombus-shaped phenocrysts of feldspar (commonly anorthoclase) embedded in a very fine-grained red-brown matrix. The composition of rhomb porphyry places it in the trachyte–latite classification of the QAPF diagram.
Rhomb porphyry is found in continental rift areas, including the East African Rift (including Mount Kilimanjaro), Mount Erebus near the Ross Sea in Antarctica, the Oslo graben in Norway, and south-central British Columbia.
Use in art and architecture
Antiquity and Byzantium
To the Romans it was known as Lapis porphyrites. Pliny the Elder's Natural History (36, 11) affirmed that the "Imperial Porphyry" had been discovered in Egypt during the reign of Tiberius; an inscription recently discovered and dated from AD 18 mentions the Roman Caius Cominius Leugas as the finder of this new quarry.
Ancient Egyptians used other decorative porphyritic stones of a very close composition and appearance, but apparently remained unaware of the presence of the Roman grade although it was located in their own country. It was also sometimes used in Minoan art, and as early as 1850 BC on Crete in Minoan Knossos there were large column bases made of porphyry.
It was called "Imperial" as the mines, as elsewhere in the empire, were owned by the emperor. The red porphyry all came from the Gabal Abu Dukhan quarry (or Mons Porphyrites) in the Eastern Desert of Egypt, from 600 million-year-old andesite of the Arabian-Nubian Shield. The road from the quarry westward to Qena (Roman Maximianopolis) on the Nile, which Ptolemy put on his second-century map, was first described by Strabo, and it is to this day known as the Via Porphyrites, the Porphyry Road, its track marked by the hydreumata, or watering wells that made it viable in this utterly dry landscape. It was used for all the red porphyry columns in Rome, the togas on busts of emperors, the panels in the revetment of the Pantheon, the Column of Constantine in Istanbul as well as the altars and vases and fountain basins reused in the Renaissance and dispersed as far as Kyiv.
The Romans also used "Green Porphyry" (lapis Lacedaemonius, from Greece, also known today as Serpentine), and "Black Porphyry" from the same Egyptian quarry.
After the fifth century the quarry was lost to sight for many centuries. Byzantium scholar Alexander Vasiliev suggested this was the consequence of the Council of Chalcedon in 451 and the subsequent troubles in Egypt. The scientific members of the French Expedition under Napoleon sought it in vain, and it was only when the Eastern Desert was reopened for study under Muhammad Ali that the site was rediscovered by the English Egyptologists James Burton and John Gardner Wilkinson in 1823.
Porphyry was extensively used in Byzantine imperial monuments, for example in Hagia Sophia and in the "Porphyra", the official delivery room for use of pregnant Empresses in the Great Palace of Constantinople, giving rise to the phrase "born in the purple".
Choosing porphyry as a material was a bold and specific statement for late Imperial Rome. As if it were not enough that porphyry was explicitly for imperial use, the stone's rarity set the emperors apart from their subjects as their superiors. The comparative vividness of porphyry to other stones underscored that these figures were not regular citizens, but many levels above, even gods, and worthy of the respect they expected. Porphyry made the emperors unapproachable in terms of power and nature, belonging to another world, the world of the mighty gods, present for a short time on earth.
Porphyry also stood in for the physical purple robes Roman emperors wore to show status, because of its purple colouring. Similar to porphyry, purple fabric was extremely difficult to make, as what we now call Tyrian purple required the use of rare sea snails to make the dye. The colour itself reminded the public how to behave in the presence of the emperors, with respect bordering on worship for the self-proclaimed god-kings.
Roman and late Roman imperial sarcophagi
A uniquely prestigious use of porphyry was its choice as material for imperial sarcophagi in the 4th and early 5th centuries. That tradition appears to have been started with Diocletian's porphyry sarcophagus in his mausoleum, which was destroyed when the building was repurposed as a church but of which probable fragments are at the Archaeological Museum in Split, Croatia. The oldest and best-preserved ones are now conserved at the Vatican Museums and known as the Sarcophagi of Helena and Constantina.
Nine other imperial porphyry sarcophagi were long held in the Church of the Holy Apostles in Constantinople. They were described by Constantine VII Porphyrogenitus in the De Ceremoniis (mid-10th century), who specified them to be respectively of Constantine the Great, Constantius II, Julian, Jovian, Theodosius I, Arcadius, Aelia Eudoxia, Theodosius II, and Marcian. Of these, most still exist in complete or fragmentary form, despite depredations by later Byzantine Emperors, Crusaders, and Ottoman conquerors. Four presently adorn the facade of the main building of the İstanbul Archaeology Museums, including one whose rounded shape led Alexander Vasiliev to suggest attribution to Emperor Julian on the basis of Constantine Porphyrogenitus's description. Vasiliev conjectures that the nine imperial sarcophagi, including one which carries a crux ansata or Egyptian cross, were carved in Egypt before shipment to Constantinople.
Porphyry sarcophagi in post-Roman Western Europe
The imperial porphyry sarcophagi tradition was emulated by Ostrogothic King Theodoric the Great (454-526), whose mausoleum in Ravenna still contains a porphyry tub that was used as his sarcophagus. Similarly Charles the Bald, King of West Francia and Roman Emperor, was buried at Saint-Denis in a porphyry tub which may be the same one known as "Dagobert's tub" (cuve de Dagobert), now in the Louvre.
The tomb of Peter III of Aragon, in the Monastery of Santes Creus near Tarragona, reuses a porphyry tub or alveus, which has been conjectured to be originally the sarcophagus of Late Roman Emperor Constans in his mausoleum at Centcelles, a nearby site with a well-preserved 4th-century rotunda.
In twelfth- and thirteenth-century Sicily, another group of porphyry sarcophagi were produced from the reign of Roger II onward and used for Royal and then Imperial burials, namely those of King Roger II, King William I, Emperor Henry VI, Empress Constance, and Emperor Frederick II. They are all now in the Palermo Cathedral, except William's in Monreale Cathedral. Scholar Rosa Bacile argues that they were carved by a local workshop from porphyry imported from Rome, the latter four plausibly (based on observation of their fluting) all from a single column shaft that may have been taken from the Baths of Caracalla or the Baths of Diocletian. She notes that these Sicilian porphyry sarcophagi "are the very first examples of medieval free-standing secular tombs in the West, and therefore play a unique role within the history of Italian sepulchral art (earlier and later tombs are adjacent to, and dependent on walls)."
Six grand porphyry sarcophagi are featured along the walls of the octagonal Cappella dei Principi (Chapel of the Princes) that was built as one of two chapels in the architectural complex of the Basilica of San Lorenzo, in Florence, Italy, for the de' Medici family. Purple porphyry was used lavishly throughout the opulent chapel as well, with a revetment of marbles, inlaid with other colored marbles and semi-precious stone, that covers the walls completely. Envisioned by Cosimo I, Grand Duke of Tuscany (1537–1574), it was initiated by Ferdinand I de' Medici, following a design by Matteo Nigetti that won an informal competition held in 1602 by Don Giovanni de' Medici (a son of Cosimo I), which was altered somewhat during execution by Buontalenti.
The tomb of Napoleon at Les Invalides in Paris, designed by architect Louis Visconti, is centered on the deceased emperor's sarcophagus that often has been described as made of red porphyry although this is incorrect. Napoleon's sarcophagus is made of quartzite, however, its pedestal is made of green andesite porphyry from Vosges. The sarcophagus of Arthur Wellesley, 1st Duke of Wellington at St Paul's Cathedral was completed in 1858. and was made from a single piece of Cornish porphyry, of a type called luxullianite, which was found in a field near Lostwithiel.
Modern uses
In countries where many automobiles have studded winter tires such as Sweden, Finland, and Norway, it is common that highways are paved with asphalt made of porphyry aggregate to make the wearing course withstand the extreme wear from the spiked winter tires.
| Physical sciences | Igneous rocks | Earth science |
305405 | https://en.wikipedia.org/wiki/Extrusive%20rock | Extrusive rock | Extrusive rock refers to the mode of igneous volcanic rock formation in which hot magma from inside the Earth flows out (extrudes) onto the surface as lava or explodes violently into the atmosphere to fall back as pyroclastics or tuff. In contrast, intrusive rock refers to rocks formed by magma which cools below the surface.
The main effect of extrusion is that the magma can cool much more quickly in the open air or under seawater, and there is little time for the growth of crystals. Sometimes, a residual portion of the matrix fails to crystallize at all, instead becoming a natural glass like obsidian.
If the magma contains abundant volatile components which are released as free gas, then it may cool with large or small vesicles (bubble-shaped cavities) such as in pumice, scoria, or vesicular basalt. Other examples of extrusive rocks are rhyolite and andesite.
Texture
The texture of extrusive rocks is characterized by fine-grained crystals indistinguishable to the human eye, described as aphantic. Crystals in aphantic rocks are small in size due to their rapid formation during eruption. Any larger crystals visible to the human eye, called phenocrysts, form earlier while slowly cooling in the magma reservoir. When igneous rocks contain two distinct grain sizes, the texture is porphyritic, and the finer crystals are called the groundmass. The extrusive rocks scoria and pumice have a vesicular, bubble-like, texture due to the presence of vapor bubbles trapped in the magma.
Extrusive bodies and rock types
Shield volcanoes are large, slow forming volcanoes that erupt fluid basaltic magma that cools to form the extrusive rock basalt. Basalt is composed of minerals readily available in the planet's crust, including feldspars and pyroxenes.
Fissure volcanoes pour out low viscosity basaltic magma from fissure vents to form the extrusive rock basalt.
Composite or stratovolcanoes often have andesitic magma and typically form the extrusive rock andesite. Andesitic magma is composed of many gases and melted mantle rocks.
Cinder or scoria cones violently expel lava with high gas content, and due to the vapor bubbles in this mafic lava, the extrusive basalt scoria is formed.
Lava domes are formed by high viscosity lava that piles up, forming a dome shape. Domes typically solidify to form the rich in silica extrusive rock obsidian and sometimes dacite domes form the extrusive rock dacite, like in the case of Mount St. Helens.
Calderas are volcanic depressions formed after an erupted volcano collapses. Resurgent calderas can refill with an eruption of rhyolitic magma to form the extrusive rock rhyolite like the Yellowstone Caldera.
Submarine volcanoes erupt on the ocean floor and produce the extrusive rock pumice. Pumice is a light-weight glass with a vesicular texture that differs from scoria in its silicic composition and therefore floats.
| Physical sciences | Igneous rocks | Earth science |
305413 | https://en.wikipedia.org/wiki/Amphibolite | Amphibolite | Amphibolite () is a metamorphic rock that contains amphibole, especially hornblende and actinolite, as well as plagioclase feldspar, but with little or no quartz. It is typically dark-colored and dense, with a weakly foliated or schistose (flaky) structure. The small flakes of black and white in the rock often give it a salt-and-pepper appearance.
Amphibolite frequently forms by metamorphism of mafic igneous rocks, such as basalt. However, because metamorphism creates minerals entirely based upon the chemistry of the protolith, certain 'dirty marls' and volcanic sediments may also metamorphose to an amphibolite assemblage. Deposits containing dolomite and siderite also readily yield amphibolite (tremolite-schist, grunerite-schist, and others) especially where there has been a certain amount of contact metamorphism by adjacent granitic masses. Metamorphosed basalt (metabasalt) creates ortho-amphibolite and other chemically appropriate lithologies create para-amphibolite.
Although tremolite is a metamorphic amphibole, it is most commonly derived from highly metamorphosed ultramafic rocks, and thus tremolite-talc schist is not generally considered a variety of amphibolite. A holocrystalline plutonic igneous rock composed primarily of hornblende amphibole is called a hornblendite, which is usually a crystal cumulate rock. Igneous rocks with greater than 90% amphiboles, which have a feldspar groundmass, may be lamprophyres.
Ortho-amphibolite vs. para-amphibolite
Metamorphic rocks composed primarily of amphibole, plagioclase, with subordinate epidote, zoisite, chlorite, quartz, titanite, and accessory leucoxene, ilmenite and magnetite which have a protolith of an igneous rock are known as ortho-amphibolite.
Para-amphibolite will generally have the same equilibrium mineral assemblage as ortho-amphibolite, with more biotite, and may include more quartz, plagioclase, and depending on the protolith, more calcite/aragonite and wollastonite.
Often the easiest way to determine the true nature of an amphibolite is to inspect its field relationships; especially whether it is interfingered with other metasedimentary rocks, especially greywacke and other poorly sorted sedimentary rocks. If the amphibolite appears to transgress apparent protolith bedding surfaces it is an ortho-amphibolite, as this suggests it was a dyke. Picking a sill and thin metamorphosed lava flows may be more troublesome.
Thereafter, whole rock geochemistry will suitably identify ortho- from para-amphibolite.
The word metabasalt was thus coined, largely to avoid the confusion between ortho-amphibolite and para-amphibolite. This term is recommended by the British Geological Survey when it is possible to determine the origin of the rock from its characteristics alone (and not from field relationships), particularly when the degree of metamorphism is low.
Amphibolite facies
Amphibolite as a rock defines a particular set of temperature and pressure conditions known as the amphibolite facies. However, caution must be applied here before embarking on metamorphic mapping based on amphibolite alone.
First, for an ortho-amphibolite or amphibolite to be classed as a metamorphic amphibolite, it must be certain that the amphibole in the rock is a prograde metamorphic product, and not a retrograde metamorphic product. For instance, actinolite amphibole is a common product of retrograde metamorphism of metabasalt at (upper) greenschist facies conditions. Often, this will take on the crystal form and habit of the original protolith assemblage; actinolite pseudomorphically replacing pyroxene is an indication that the amphibolite may not represent a peak metamorphic grade in the amphibolite facies.
Actinolite schist is often the result of hydrothermal alteration or metasomatism, and thus may not, necessarily, be a good indicator of metamorphic conditions when taken in isolation.
Second, the microstructure and crystal size of the rock must be appropriate. Amphibolite facies conditions are experienced at temperatures in excess of 500 °C and pressures less than 1.2 GPa, well within the ductile deformation field. Gneissic texture may occur nearby, if not then mylonite zones, foliations and ductile behaviour, including stretching lineations may occur.
While it is not impossible to have remnant protolith mineralogy, this is rare. More common is to find phenocrysts of pyroxene, olivine, plagioclase and even magmatic amphibole such as pargasite rhombohedra, pseudomorphed by hornblende amphibole. Original magmatic textures, especially crude magmatic layering in layered intrusions, is often preserved.
Amphibolite facies equilibrium mineral assemblages of various protolith rock types consist of:
Basalt ortho-amphibolite; hornblende/actinolite +/- albite +/- biotite +/- quartz +/- accessories; often remnant greenschist facies assemblages including, notably, chlorite
High-magnesia basalt; as ortho-amphibolite, but may contain anthophyllite, a Mg-rich amphibole
Ultramafic rocks; tremolite, asbestiform amphibole, talc, pyroxene, wollastonite, prograde metamorphic olivine (rarely)
Sedimentary para-amphibolite; hornblende/actinolite +/- albite +/- biotite +/- quartz +/- garnet (calcite +/- wollastonite)
Pelite; quartz, orthoclase +/- albite, +/- biotite +/- actinolite +/- garnet +/- staurolite +/- sillimanite
Amphibolite facies is usually a product of Barrovian Facies Sequence or advanced Abukuma Facies Sequence metamorphic trajectories. Amphibolite facies is a result of continuing burial and thermal heating after greenschist facies is exceeded. Further burial and metamorphic compression (but little extra heat) will lead to eclogite facies metamorphism; with more advanced heating the majority of rocks begin melting in excess of 650 to 700 °C in the presence of water. In dry rocks, however, additional heat (and burial) may result in granulite facies conditions.
Uralite
Uralite is a particular hydrothermally altered pyroxenite; during autogenic hydrothermal circulation the primary mineralogy of pyroxene and plagioclase, etc. has altered to actinolite and saussurite (albite + epidote). The texture is distinctive, the pyroxene altered to fuzzy, radially arranged actinolite pseudomorphically after pyroxene, and saussuritised plagioclase.
Epidiorite
The archaic term epidiorite is sometimes used, especially in Europe, to refer to a metamorphosed ortho-amphibolite with a protolith of diorite, gabbro or other mafic intrusive rock. In epidiorite the original clinopyroxene (most often augite) has been replaced by the fibrous amphibole uralite.
Uses
Amphibolite was a favourite material for the production of adzes (shoe-last-celts) in the central European early Neolithic (Linearbandkeramic and Rössen cultures).
Amphibolite is a common dimension stone used in construction, paving, facing of buildings, especially because of its attractive textures, dark color, hardness and polishability and its ready availability.
| Physical sciences | Metamorphic rocks | Earth science |
305429 | https://en.wikipedia.org/wiki/Feldspathoid | Feldspathoid | The feldspathoids are a group of tectosilicate minerals which resemble feldspars but have a different structure and much lower silica content. They occur in rare and unusual types of igneous rocks, and are usually not found in rocks containing primary quartz. A notable exception where feldspathoids and quartz-bearing rocks are found together is the Red Hill Syenite.
Foid, a contraction of the term feldspathoid, is applied to any igneous rock containing up to 60% modal feldspathoid minerals. For example, a syenite with significant nepheline present can be termed a nepheline-bearing syenite or nepheline syenite, with the term nepheline replaceable by any foid mineral. Such terminology is used in the Streckeisen (QAPF) classification of igneous rocks.
Feldspathoid minerals
Sodalite Group
| Physical sciences | Silicate minerals | Earth science |
305465 | https://en.wikipedia.org/wiki/Lens%20%28vertebrate%20anatomy%29 | Lens (vertebrate anatomy) | The lens, or crystalline lens, is a transparent biconvex structure in most land vertebrate eyes. Relatively long, thin fiber cells make up the majority of the lens. These cells vary in architecture and are arranged in concentric layers. New layers of cells are recruited from a thin epithelium at the front of the lens, just below the basement membrane surrounding the lens. As a result the vertebrate lens grows throughout life. The surrounding lens membrane referred to as the lens capsule also grows in a systematic way, ensuring the lens maintains an optically suitable shape in concert with the underlying fiber cells. Thousands of suspensory ligaments are embedded into the capsule at its largest diameter which suspend the lens within the eye. Most of these lens structures are derived from the epithelium of the embryo before birth.
Along with the cornea, aqueous, and vitreous humours, the lens refracts light, focusing it onto the retina. In many land animals the shape of the lens can be altered, effectively changing the focal length of the eye, enabling them to focus on objects at various distances. This adjustment of the lens is known as accommodation (see also below). In many fully aquatic vertebrates, such as fish, other methods of accommodation are used, such as changing the lens's position relative to the retina rather than changing the shape of the lens. Accommodation is analogous to the focusing of a photographic camera via changing its lenses. In land vertebrates the lens is flatter on its anterior side than on its posterior side, while in fish the lens is often close to spherical.
Accommodation in humans is well studied and allows artificial means of supplementing our focus, such as glasses, for correction of sight as we age. The refractive power of a younger human lens in its natural environment is approximately 18 dioptres, roughly one-third of the eye's total power of about 60 dioptres. By 25 years of age the ability of the lens to alter the light path has reduced to 10 dioptres and accommodation continues to decline with age.
Structure
Position in the eye
The lens is located towards the front part of the vertebrate eye, called the anterior segment, which includes the cornea and iris positioned in front of the lens. The lens is held in place by the suspensory ligaments (Zonule of Zinn), attaching the lens at its equator to the rest of the eye through the ciliary body. Behind the lens is the jelly-like vitreous body which helps hold the lens in place. At the front of the lens is the liquid aqueous humor which bathes the lens with nutrients and other things. Land vertebrate lenses usually have an ellipsoid, biconvex shape. The front surface is less curved than the back. In a human adult, the lens is typically about 10mm in diameter and 4mm thick, though its shape changes with accommodation and its size grows throughout a person's lifetime.
Anatomy
The lens has three main parts: the lens capsule, the lens epithelium, and the lens fibers. The lens capsule is a relatively thick basement membrane forming the outermost layer of the lens. Inside the capsule, much thinner lens fibers form the bulk of the lens. The cells of the lens epithelium form a thin layer between the lens capsule and the outermost layer of lens fibers at the front of the lens but not the back. The lens itself lacks nerves, blood vessels, or connective tissue. Anatomists will often refer to positions of structures in the lens by describing it like a globe of the world. The front and back of the lens are referred to as the anterior and posterior "poles", like the North and South poles. The "equator" is the outer edge of the lens often hidden by the iris and is the area of most cell differentiation. As the equator is not generally in the light path of the eye, the structures involved with metabolic activity avoid scattering light that would otherwise affect vision.
Lens capsule
The lens capsule is a smooth, transparent basement membrane that completely surrounds the lens. The capsule is elastic and its main structural component is collagen. It is presumed to be synthesized by the lens epithelium and its main components in order of abundance are heparan sulfate proteoglycan (sulfated glycosaminoglycans (GAGs)), entactin, type IV collagen and laminin. The capsule is very elastic and so allows the lens to assume a more spherical shape when the tension of the suspensory ligaments is reduced. The human capsule varies from 2 to 28 micrometres in thickness, being thickest near the equator (peri-equatorial region) and generally thinner near the posterior pole.
The photos from electron and light microscopes show an area of the capsule lens equator where the capsule grows and adjacent to where thousands of suspensory ligaments attach. Attachment must be strong enough to stop the ligaments being detached from the lens capsule. Forces are generated from holding the lens in place and the forces added to during focusing. While the capsule is thinnest at the equator where its area is increasing, the anterior and posterior capsule is thinner than the area of ligament attachment.
Lens epithelium
The lens epithelium is a single layer of cells at the front of the lens between the lens capsule and the lens fibers. By providing the lens fibers with nutrients and removing waste, the cells of the epithelium maintain lens homeostasis. As ions, nutrients, and liquid enter the lens from the aqueous humor, Na+/K+-ATPase pumps in the lens epithelial cells pump ions out of the lens to maintain appropriate lens osmotic concentration and volume, with equatorially positioned lens epithelium cells contributing most to this current. The activity of the Na+/K+-ATPases keeps water and current flowing through the lens from the poles and exiting through the equatorial regions.
The cells of the lens epithelium also divide into new lens fibers at the lens equator. The lens lays down fibers from when it first forms in embryo until death.
Lens fibers
The lens fibers form the bulk of the lens. They are long, thin, transparent cells, firmly packed, with diameters typically 4–7 micrometres and lengths of up to 12mm long in humans. The lens fibers stretch lengthwise from the posterior to the anterior poles and, when cut horizontally, are arranged in concentric layers rather like the layers of an onion. If cut along the equator, cells have a hexagonal cross section, appearing as a honeycomb. The approximate middle of each fiber lies around the equator. These tightly packed layers of lens fibers are referred to as laminae. The lens fiber cytoplasms are linked together via gap junctions, intercellular bridges and interdigitations of the cells that resemble "ball and socket" forms.
The lens is split into regions depending on the age of the lens fibers of a particular layer. Moving outwards from the central, oldest layer, the lens is split into an embryonic nucleus, the fetal nucleus, the adult nucleus, the inner and outer cortex. New lens fibers, generated from the lens epithelium, are added to the outer cortex. Mature lens fibers have no organelles or nuclei.
Cell fusion, voids and vacuoles
With the advent of other ways of looking at cellular structures of lenses while still in the living animal it became apparent that regions of fiber cells, at least at the lens anterior, contain large voids and vacuoles. These are speculated to be involved in lens transport systems linking the surface of the lens to deeper regions. Very similar looking structures also indicate cell fusion in the lens. The cell fusion is shown by micro-injection to form a stratified syncytium in whole lens cultures.
Development
Development of the vertebrate lens begins when the human embryo is about 4mm long. The accompanying picture shows the process in a more easily studied chicken embryo. Unlike the rest of the eye which is derived mostly from the inner embryo layers, the lens is derived from the skin around the embryo. The first stage of lens formation takes place when a sphere of cells formed by budding of the inner embryo layers comes close to the embyro's outer skin. The sphere of cells induces nearby outer skin to start changing into the lens placode. The lens placode is the first stage of transformation of a patch of skin into the lens. At this early stage, the lens placode is a single layer of cells.
As development progresses, the lens placode begins to deepen and bow inwards. As the placode continues to deepen, the opening to the surface ectoderm constricts and the lens cells bud off from the embryo's skin to form a sphere of cells known as the "lens vesicle". When the embryo is about 10mm long the lens vesicle has completely separated from the skin of the embryo.
The embryo then sends signals from the developing retina, inducing the cells closest to the posterior end of the lens vesicle to elongate toward the anterior end of the vesicle. These signals also induce the synthesis of proteins called crystallins. As the name suggests the crystallins can form a clear highly refractive jelly. These elongating cells eventually fill in the center of the vesicle with cells, that are long and thin like a strand of hair, called fibers. These primary fibers become the nucleus in the mature lens. The epithelial cells that do not form into fibers nearest the lens front give rise to the lens epithelium.
Additional fibers are derived from lens epithelial cells located at the lens equator. These cells lengthen towards the front and back wrapping around fibers already laid down. The new fibers need to be longer to cover earlier fibers but as the lens gets larger the ends of the newer fibers no longer reach as far towards the front and back of the lens. The lens fibers that do not reach the poles form tight, interdigitating seams with neighboring fibers. These seams being less crystalline than the bulk of the lens are more visible and are termed "sutures". The suture patterns become more complex as more layers of lens fibers are added to the outer portion of the lens.
The lens continues to grow after birth, with the new secondary fibers being added as outer layers. New lens fibers are generated from the equatorial cells of the lens epithelium, in a region referred to as the "germinative zone" and "bow region". The lens epithelial cells elongate, lose contact with the capsule and epithelium at the back and front of the lens, synthesize crystallin, and then finally lose their nuclei (enucleate) as they become mature lens fibers. In humans, as the lens grows by laying down more fibers through to early adulthood, the lens becomes more ellipsoid in shape. After about age 20 the lens grows rounder again and the iris is very important for this development.
Several proteins control the embryonic development of the lens though PAX6 is considered the master regulator gene of this organ. Other effectors of proper lens development include the Wnt signaling components BCL9 and Pygo2. The whole process of differentiation of the epithelial cells into crystallin filled fiber cells without organelles occurs within the confines of the lens capsule. Older cells cannot be shed and are instead internalized towards the center of the lens. This process results in a complete temporally layered record of the differentiation process from the start at the lens surface to the end at the lens center. The lens is therefore valuable to scientists studying the process of cell differentiation.
Variations in lens structure
In many aquatic vertebrates, the lens is considerably thicker, almost spherical resulting in increased light refraction. This difference helps compensate for the smaller angle of refraction between the eye's cornea and the watery environment, as they have more similar refractive indices than cornea and air. The fiber cells of fish are generally considerably thinner than those of land vertebrates and it appears crystallin proteins are transported to the organelle free cells at the lens exterior to the inner cells through many layers of cells. Some vertebrates need to see well both above and below water at times. One example is diving birds which have the ability to change focus by 50 to 80 dioptres. Compared with animals adapted for only one environment diving birds have a somewhat altered lens and cornea structure with focus mechanisms to allow for both environments. Even among terrestrial animals the lens of primates such as humans is unusually flat going some way to explain why our vision, unlike diving birds, is particularly blurry under water.
Function
Focusing
In humans the widely quoted Helmholtz mechanism of focusing, also called accommodation, is often referred to as a "model". Direct experimental proof of any lens model is necessarily difficult as the vertebrate lens is transparent and only functions well in the living animals. When considering all vertebrates aspects of all models may play varying roles in lens focus.
The shape changing lens of many land based vertebrates
External forces
The model of a shape changing lens of humans was proposed by Young in a lecture on the 27th Nov 1800. Others such as Helmholtz and Huxley refined the model in the mid-1800s explaining how the ciliary muscle contracts rounding the lens to focus near and this model was popularized by Helmholtz in 1909. The model may be summarized like this. Normally the lens is held under tension by its suspending ligaments being pulled tight by the pressure of the eyeball. At short focal distance the ciliary muscle contracts relieving some of the tension on the ligaments, allowing the lens to elastically round up a bit, increasing refractive power. Changing focus to an object at a greater distance requires a thinner less curved lens. This is achieved by relaxing some of the sphincter like ciliary muscles. While not referenced this presumably allows the pressure in the eyeball to again expand it outwards, pulling harder on the lens making it less curved and thinner, so increasing the focal distance. There is a problem with the Helmholtz model in that despite mathematical models being tried none has come close enough to working using only the Helmholtz mechanisms.
Schachar has proposed a model for land based vertebrates that was not well received. The theory allows mathematical modeling to more accurately reflect the way the lens focuses while also taking into account the complexities in the suspensory ligaments and the presence of radial as well as circular muscles in the ciliary body. In this model the ligaments may pull to varying degrees on the lens at the equator using the radial muscles while the ligaments offset from the equator to the front and back are relaxed to varying degrees by contracting the circular muscles. These multiple actions operating on the elastic lens allows it to change lens shape at the front more subtly. Not only changing focus, but also correcting for lens aberrations that might otherwise result from the changing shape while better fitting mathematical modeling.
The "catenary" model of lens focus proposed by Coleman demands less tension on the ligaments suspending the lens. Rather than the lens as a whole being stretched thinner for distance vision and allowed to relax for near focus, contraction of the circular ciliary muscles results in the lens having less hydrostatic pressure against its front. The lens front can then reform its shape between the suspensory ligaments in a similar way to a slack chain hanging between two poles might change its curve when the poles are moved closer together. This model requires fluid movement of the lens front only rather than trying to change the shape of the lens as a whole.
Internal forces
When Thomas Young proposed the changing of the human lens's shape as the mechanism for focal accommodation in 1801 he thought the lens may be a muscle capable of contraction. This type of model is termed intracapsular accommodation as it relies on activity within the lens. In a 1911 Nobel lecture Allvar Gullstrand spoke on "How I found the intracapsular mechanism of accommodation" and this aspect of lens focusing continues to be investigated. Young spent time searching for the nerves that could stimulate the lens to contract without success. Since that time it has become clear the lens is not a simple muscle stimulated by a nerve so the 1909 Helmholtz model took precedence. Pre-twentieth century investigators did not have the benefit of many later discoveries and techniques. Membrane proteins such as aquaporins which allow water to flow into and out of cells are the most abundant membrane protein in the lens. Connexins which allow electrical coupling of cells are also prevalent. Electron microscopy and immunofluorescent microscopy show fiber cells to be highly variable in structure and composition. Magnetic resonance imaging confirms a layering in the lens that may allow for different refractive plans within it. The refractive index of human lens varies from approximately 1.406 in the central layers down to 1.386 in less dense layers of the lens. This index gradient enhances the optical power of the lens. As more is learned about mammalian lens structure from in situ Scheimpflug photography, MRI and physiological investigations it is becoming apparent the lens itself is not responding entirely passively to the surrounding ciliary muscle but may be able to change its overall refractive index through mechanisms involving water dynamics in the lens still to be clarified. The accompanying micrograph shows wrinkled fibers from a relaxed sheep lens after it is removed from the animal indicating shortening of the lens fibers during near focus accommodation. The age related changes in the human lens may also be related to changes in the water dynamics in the lens.
Lenses of birds, reptiles, amphibians, fish and others
In reptiles and birds, the ciliary body which supports the lens via suspensory ligaments also touches the lens with a number of pads on its inner surface. These pads compress and release the lens to modify its shape while focusing on objects at different distances; the suspensory ligaments usually perform this function in mammals. With vision in fish and amphibians, the lens is fixed in shape, and focusing is instead achieved by moving the lens forwards or backwards within the eye using a muscle called the retractor lentus.
In cartilaginous fish, the suspensory ligaments are replaced by a membrane, including a small muscle at the underside of the lens. This muscle pulls the lens forward from its relaxed position when focusing on nearby objects. In teleosts, by contrast, a muscle projects from a vascular structure in the floor of the eye, called the falciform process, and serves to pull the lens backwards from the relaxed position to focus on distant objects. While amphibians move the lens forward, as do cartilaginous fish, the muscles involved are not similar in either type of animal. In frogs, there are two muscles, one above and one below the lens, while other amphibians have only the lower muscle.
In the simplest vertebrates, the lampreys and hagfish, the lens is not attached to the outer surface of the eyeball at all. There is no aqueous humor in these fish, and the vitreous body simply presses the lens against the surface of the cornea. To focus its eyes, a lamprey flattens the cornea using muscles outside of the eye and pushes the lens backwards.
While not vertebrate, brief mention is made here of the convergent evolution of vertebrate and Molluscan eyes. The most complex Molluscan eye is the Cephalopod eye which is superficially similar structure and function to a vertebrate eye, including accommodation, while differing in basic ways such as having a two part lens and no cornea. The fundamental requirements of optics must be filled by all eyes with lenses using the tissues at their disposal so superficially eyes all tend to look similar. It is the way optical requirements are met using different cell types and structural mechanisms that varies among animals.
Crystallins and transparency
Crystallins are water-soluble proteins that compose over 90% of the protein within the lens. The three main crystallin types found in the human eye are α-, β-, and γ-crystallins. Crystallins tend to form soluble, high-molecular weight aggregates that pack tightly in lens fibers, thus increasing the index of refraction of the lens while maintaining its transparency. β and γ crystallins are found primarily in the lens, while subunits of α -crystallin have been isolated from other parts of the eye and the body. α-crystallin proteins belong to a larger superfamily of molecular chaperone proteins, and so it is believed that the crystallin proteins were evolutionarily recruited from chaperone proteins for optical purposes. The chaperone functions of α-crystallin may also help maintain the lens proteins, which must last a human for their entire lifetime.
Another important factor in maintaining the transparency of the lens is the absence of light-scattering organelles such as the nucleus, endoplasmic reticulum, and mitochondria within the mature lens fibers. Lens fibers also have a very extensive cytoskeleton that maintains the precise shape and packing of the lens fibers; disruptions/mutations in certain cytoskeletal elements can lead to the loss of transparency.
The lens blocks most ultraviolet light in the wavelength range of 300–400 nm; shorter wavelengths are blocked by the cornea. The pigment responsible for blocking the light is 3-hydroxykynurenine glucoside, a product of tryptophan catabolism in the lens epithelium. High intensity ultraviolet light can harm the retina, and artificial intraocular lenses are therefore manufactured to also block ultraviolet light. People lacking a lens (a condition known as aphakia) perceive ultraviolet light as whitish blue or whitish-violet.
Nourishment
The lens is metabolically active and requires nourishment in order to maintain its growth and transparency. Compared to other tissues in the eye, however, the lens has considerably lower energy demands.
By nine weeks into human development, the lens is surrounded and nourished by a net of vessels, the tunica vasculosa lentis, which is derived from the hyaloid artery. Beginning in the fourth month of development, the hyaloid artery and its related vasculature begin to atrophy and completely disappear by birth. In the postnatal eye, Cloquet's canal marks the former location of the hyaloid artery.
After regression of the hyaloid artery, the lens receives all its nourishment from the aqueous humor. Nutrients diffuse in and waste diffuses out through a constant flow of fluid from the anterior/posterior poles of the lens and out of the equatorial regions, a dynamic that is maintained by the Na+/K+-ATPase pumps located in the equatorially positioned cells of the lens epithelium. The interaction of these pumps with water channels into cells called aquaporins, molecules less than 100 daltons in size among cells via gap junctions, and calcium using transporters/regulators (TRPV channels) results in a flow of nutrients throughout the lens.
Glucose is the primary energy source for the lens. As mature lens fibers do not have mitochondria, approximately 80% of the glucose is metabolized via anaerobic metabolism. The remaining fraction of glucose is shunted primarily down the pentose phosphate pathway. The lack of aerobic respiration means that the lens consumes very little oxygen.
Clinical significance
Cataracts are opacities of the lens. While some are small and do not require any treatment, others may be large enough to block light and obstruct vision. Cataracts usually develop as the aging lens becomes more and more opaque, but cataracts can also form congenitally or after injury to the lens. Nuclear sclerosis is a type of age-related cataract. Diabetes is another risk factor for cataract. Cataract surgery involves the removal of the lens and insertion of an artificial intraocular lens.
Presbyopia is the age-related loss of accommodation, which is marked by the inability of the eye to focus on nearby objects. The exact mechanism is still unknown, but age-related changes in the hardness, shape, and size of the lens have all been linked to the condition.
Ectopia lentis is the displacement of the lens from its normal position.
Aphakia is the absence of the lens from the eye. Aphakia can be the result of surgery or injury, or it can be congenital.
Additional images
| Biology and health sciences | Visual system | Biology |
305649 | https://en.wikipedia.org/wiki/Oceanic%20dolphin | Oceanic dolphin | Oceanic dolphins or Delphinidae are a widely distributed family of dolphins that live in the sea. Close to forty extant species are recognised. They include several big species whose common names contain "whale" rather than "dolphin", such as the Globicephalinae (round-headed whales, which include the false killer whale and pilot whale). Delphinidae is a family within the superfamily Delphinoidea, which also includes the porpoises (Phocoenidae) and the Monodontidae (beluga whale and narwhal). River dolphins are relatives of the Delphinoidea.
Oceanic dolphins range in size from the and Maui's dolphin to the and orca, the largest known dolphin. Several species exhibit sexual dimorphism; the males are larger than females. They have streamlined muscular bodies and two limbs that are modified into flippers. Though not quite as flexible as seals, some dolphins can travel at speeds 29 km/h (18 mph) for short distances. Most delphinids primarily eat fish, along with a smaller number of squid and small crustaceans, but some species specialise in eating squid, or, in the case of the orca, also eat marine mammals and birds. All, however, are purely carnivorous. They typically have between 100 and 200 teeth, although a few species have considerably fewer. Delphinids travel in large pods, which may number a thousand individuals in some species. Each pod forages over a range of tens to hundreds of square kilometres. Some pods have a loose social structure, with individuals frequently joining or leaving, but others seem to be more permanent, perhaps dominated by a male and a harem of females. Individuals communicate by sound, producing low-frequency whistles, and also produce high-frequency broadband clicks of 80–220 kHz, which are primarily used for echolocation. Gestation lasts from 10 to 18 months, and results in the birth of a single calf. Some species are well adapted for diving to great depths. They have a layer of fat, or blubber, under the skin to keep warm in the cold water.
Although oceanic dolphins are widespread, most species prefer the warmer waters of the tropic zones, but some, like the right whale dolphin, prefer colder climates. Some have a global distribution, like the orca. Oceanic dolphins feed largely on fish and squid, but a few, like the orca, feed on large mammals, like seals. Male dolphins typically mate with multiple females every year, but females only mate every two to three years. Calves are typically born in the spring and summer, and females bear all the responsibility for raising them. Mothers of some species fast and nurse their young for relatively long times. Dolphins produce a variety of vocalizations, usually in the form of clicks and whistles.
Oceanic dolphins are sometimes hunted in places such as Japan, in an activity known as dolphin drive hunting. Besides drive hunting, they also face threats from bycatch, habitat loss, and marine pollution. Dolphins have been depicted in various cultures worldwide. They occasionally feature in literature and film, as in the Warner Bros film Free Willy. Dolphins are sometimes kept in captivity and trained to perform in shows. The most common species of dolphin in captivity is the bottlenose dolphin, and less than 50 orca were found in oceanariums in 2012.
Taxonomy
Delphinids, especially bottlenose dolphins, are able to hybridize with a wide variety of other delphinid species; wholphins are just one of many possible hybrids.
Six species, sometimes referred to as "blackfish", are dolphins commonly called whales: the orca, the melon-headed whale, the pygmy killer whale, the false killer whale, and the two species of pilot whales, but they are classified under the family Delphinidae.
Recent molecular analyses indicate that several delphinid genera (especially Stenella and Lagenorhynchus) are not monophyletic as currently recognized. Thus, significant taxonomic revisions within the family are likely.
Extinct taxon, according to the Paleobiology Database.
Subfamily Delphininae
†Astadelphis
†Etruridelphis
†Septidelphis
Subfamily Globicephalinae
†Platalearostrum
†Protoglobicephala
†Rododelphis
†Arimidelphis
†Australodelphis
†Eodelphinus
†Hemisyntrachelus
†Norisdelphis
†Pliodelphis
Biology
Anatomy
The Delphinidae are the most diverse of the cetacean families, with numerous variations between species. They range in size from and (Haviside's dolphin), to and 10 tonnes (orca). Most species weigh between . They typically have curved dorsal fins, clear 'beaks' at the front of their heads, and forehead melons, although exceptions to all of these rules are found. They have a wide range of colors and patterns.
Oceanic dolphins have a torpedo-shaped, muscular body with an inflexible neck, limbs modified into flippers, nonexistent external ear flaps, a large tail fin, and a bulbous head. A dolphin skull has small eye orbits, a long snout, but not as long as its river dolphin counterpart, and eyes placed on the sides of its head. Several species exhibit sexual dimorphism, with the males being larger than the females.
Breathing involves expelling stale air from the blowhole, forming an upward, steamy spout, followed by inhaling fresh air into the lungs; a spout only occurs when the warm air from the lungs meets the cold external air, so it may only form in colder climates.
All oceanic dolphins have a thick layer of blubber, the thickness of which depends on how far the species lives from the equator. This blubber also helps keep the animal warm providing insulation from the harsh climate or cold depths. It can also aid in protection to some extent as predators would have a hard time getting through a thick layer of fat. Calves are born with only a thin layer of blubber, but some species compensate for this with lanugos.
Locomotion
Oceanic dolphins have two flippers on the underside toward the head, a dorsal fin, and a tail fin. These flippers contain four digits. Although oceanic dolphins do not possess fully developed hind limbs, some possess discrete rudimentary appendages, which may contain feet and digits. Oceanic dolphins are fast swimmers in comparison to seals who typically cruise at ; the orca, in comparison, can travel at speeds of up to . The fusing of the neck vertebrae, while increasing stability when swimming at high speeds, decreases flexibility, which means they are unable to turn their heads. Oceanic dolphins swim by moving their tail fin and rear body vertically, while their flippers are mainly used for steering. Some species log out of the water, which may allow them to travel faster. Their skeletal anatomy allows them to be fast swimmers. All species have a dorsal fin to prevent themselves from involuntarily spinning in the water.
Senses
The oceanic dolphin ear is specifically adapted to the marine environment. In humans, the middle ear works as an impedance equalizer between the outside air's low impedance and the cochlear fluid's high impedance. In dolphins, and other marine mammals, there is no great difference between the outer and inner environments. Instead of sound passing through the outer ear to the middle ear, dolphins receive sound through the throat, from which it passes through a low-impedance fat-filled cavity to the inner ear. The dolphin ear is acoustically isolated from the skull by air-filled sinus pockets, which allow for greater directional hearing underwater. Dolphins send out high frequency clicks from an organ known as a melon. This melon consists of fat, and the skull of any such creature containing a melon will have a large depression. This allows dolphins to produce biosonar for orientation. Though most dolphins do not have hair, they do have hair follicles that may perform some sensory function. Beyond locating an object, echolocation also provides the animal with an idea on an object's shape and size, though how exactly this works is not yet understood.
The eye of oceanic dolphins is relatively small for its size, yet they do retain a good degree of eyesight. As well as this, the eyes of a dolphin are placed on the sides of its head, so their vision consists of two fields, rather than a binocular view like humans have. When dolphins surface, their lens and cornea correct the nearsightedness that results from the refraction of light; they contain both rod and cone cells, meaning they can see in both dim and bright light. Dolphins do, however, lack short wavelength sensitive visual pigments in their cone cells indicating a more limited capacity for color vision than most mammals. Most dolphins have slightly flattened eyeballs, enlarged pupils (which shrink as they surface to prevent damage), slightly flattened corneas and a tapetum lucidum; these adaptations allow for large amounts of light to pass through the eye and, therefore, a very clear image of the surrounding area. They also have glands on the eyelids and outer corneal layer that act as protection for the cornea.
The olfactory lobes are absent in oceanic dolphins, suggesting that they have no sense of smell.
Oceanic dolphins are not thought to have a good sense of taste, as their taste buds are atrophied or missing altogether. However, some have preferences between different kinds of fish, indicating some sort of attachment to taste.
Behavior
Feeding
Most delphinids primarily eat fish, along with a smaller number of squid and small crustaceans, but some species specialize in eating squid, or, in the case of the orca, also eat marine mammals and birds. All, however, are purely carnivorous. They typically have between 100 and 200 teeth, although a few species have considerably fewer. Various methods of feeding exist among and within oceanic species, some apparently exclusive to a single population. Fish and squid are the main food, but the false killer whale and the orca also feed on other marine mammals. Orca on occasion also hunt whales larger than themselves.
One common feeding method is herding, where a pod squeezes a school of fish into a bait ball. Individual members then take turns plowing through the ball, feeding on the stunned fish. Corralling is a method where dolphins chase fish into shallow water to catch them more easily. Orca and bottlenose dolphins have also been known to drive their prey onto a beach to feed on it, a behavior known as beach or strand feeding. Some species also whack fish with their flukes, stunning them and sometimes knocking them out of the water.
Vocalizations
Oceanic dolphins are capable of making a broad range of sounds using nasal airsacs located just below the blowhole. Roughly three categories of sounds can be identified: frequency modulated whistles, burst-pulsed sounds, and clicks. Dolphins communicate with whistle-like sounds produced by vibrating connective tissue, similar to the way human vocal cords function, and through burst-pulsed sounds, though the nature and extent of that ability is not known. The clicks are directional and are for echolocation, often occurring in a short series called a click train. The click rate increases when approaching an object of interest. Dolphin echolocation clicks are amongst the loudest sounds made by marine animals.
Bottlenose dolphins have been found to have signature whistles. These whistles are used in order for dolphins to communicate with one another by identifying an individual. It can be seen as the dolphin equivalent of a name for humans. These signature whistles are developed during a dolphin's first year; it continues to maintain the same sound throughout its lifetime. In order to obtain each individual whistle sound, dolphins undergo vocal production learning. This consists of an experience with other dolphins that modifies the signal structure of an existing whistle sound. An auditory experience influences the whistle development of each dolphin. Dolphins are able to communicate to one another by addressing another dolphin through mimicking their whistle. The signature whistle of a male bottlenose dolphin tends to be similar to that of his mother, while the signature whistle of a female bottlenose dolphin tends to be more distinguishable. Bottlenose dolphins have a strong memory when it comes to these signature whistles, as they are able to relate to a signature whistle of an individual they have not encountered for over twenty years. Research done on signature whistle usage by other dolphin species is relatively limited. The research on other species done so far has yielded varied outcomes and inconclusive results.
Surfacing behavior and play
Oceanic dolphins frequently leap above the water surface, this being done for various reasons. When travelling, jumping can save the dolphin energy as there is less friction while in the air. This type of travel is known as porpoising. Other reasons include orientation, social displays, fighting, non-verbal communication, entertainment and attempting to dislodge parasites.
Dolphins show various types of playful behavior, often including objects, self-made bubble rings, other dolphins, or other animals. When playing with objects or small animals, common behavior includes carrying the object or animal along using various parts of the body, passing it along to other members of the group or taking it from another member, or throwing it out of the water. Dolphins have also been observed harassing animals in other ways, for example by dragging birds underwater without showing any intent to eat them. Playful behaviour that involves another animal species with active participation of the other animal can also be observed, however. Playful human interaction with dolphins is one example, but playful interactions have been observed in the wild with a number of other species as well, including humpback whales and dogs.
Intelligence
Oceanic dolphins are known to teach, learn, cooperate, scheme, and grieve. The neocortex of many species is home to elongated spindle neurons that, prior to 2007, were known only in hominids. In humans, these cells are involved in social conduct, emotions, judgment, and theory of mind. Cetacean spindle neurons are found in areas of the brain that are homologous to where they are found in humans, suggesting that they perform a similar function.
Brain size was previously considered a major indicator of the intelligence of an animal. Since most of the brain is used for maintaining bodily functions, greater ratios of brain to body mass may increase the amount of brain mass available for more complex cognitive tasks. Allometric analysis indicates that mammalian brain size scales at approximately the ⅔ or ¾ exponent of the body mass. Comparison of a particular animal's brain size with the expected brain size based on such allometric analysis provides an encephalization quotient that can be used as another indication of animal intelligence. Orca have the second largest brain mass of any animal on earth, next to the sperm whale. The brain to body mass ratio in some is second only to humans.
Self-awareness is seen, by some, to be a sign of highly developed, abstract thinking. Self-awareness, though not well-defined scientifically, is believed to be the precursor to more advanced processes like meta-cognitive reasoning (thinking about thinking) that are typical of humans. Research in this field has suggested that cetaceans, among others, possess self-awareness.
The most widely used test for self-awareness in animals is the mirror test in which a temporary dye is placed on an animal's body, and the animal is then presented with a mirror; they then see if the animal shows signs of self-recognition.
In 1995, Marten and Psarakos used television to test dolphin self-awareness. They showed dolphins real-time footage of themselves, recorded footage, and another dolphin. They concluded that their evidence suggested self-awareness rather than social behavior. While this particular study has not been repeated since then, dolphins have since passed the mirror test.
Interactions with humans
Threats
Consumption
In some parts of the world, such as Taiji, Japan and the Faroe Islands, dolphins are traditionally considered as food, and are killed in harpoon or drive hunts.
Dolphin meat is consumed in a small number of countries worldwide, which include Japan and Peru (where it is referred to as chancho marino, or "sea pork").
Dolphin meat is dense and such a dark shade of red as to appear black. Fat is located in a layer of blubber between the meat and the skin. When dolphin meat is eaten in Japan, it is often cut into thin strips and eaten raw as sashimi, garnished with onion and either horseradish or grated garlic, much as with sashimi of whale or horse meat (basashi). When cooked, dolphin meat is cut into bite-size cubes and then batter-fried or simmered in a miso sauce with vegetables. Cooked dolphin meat has a flavor very similar to beef liver. Dolphin meat is high in mercury, and may pose a health danger to humans when consumed.
There have been human health concerns associated with the consumption of dolphin meat in Japan after tests showed that dolphin meat contained high levels of mercury. There are no known cases of mercury poisoning as a result of consuming dolphin meat, though the government continues to monitor people in areas where dolphin meat consumption is high. The Japanese government recommends that children and pregnant women avoid eating dolphin meat on a regular basis.
Similar concerns exist with the consumption of dolphin meat in the Faroe Islands, where prenatal exposure to methylmercury and PCBs primarily from the consumption of pilot whale meat has resulted in neuropsychological deficits amongst children.
Legally consuming dolphin meat in the United States would be near impossible for most due to the Marine Mammal Protection Act, which forbids "...the act of hunting, killing, capture, and/or harassment of any marine mammal..." (Exceptions are made for certain groups of people, such as Alaska Natives.) Theoretically, one could only eat the meat of a dolphin which died of natural causes, which would likely be highly undesirable (and potentially dangerous).
Fishing
Various fishing methods, like seine fishing for tuna and the use of drift and gill nets, unintentionally kill many oceanic dolphins. Accidental bycatch in gill nets is common and poses a risk for mainly local dolphin populations.
Dolphin safe labels attempt to reassure consumers that fish and other marine products have been caught in a dolphin-friendly way. The earliest campaigns with "Dolphin safe" labels were initiated in the 1980s as a result of cooperation between marine activists and the major tuna companies, and involved decreasing incidental dolphin kills by up to 50% by changing the type of nets used to catch tuna. The dolphins are netted only while fishermen are in pursuit of smaller tuna. Albacore are not netted this way, making albacore the only truly dolphin-safe tuna.
Sonar
Loud underwater noises, such as those resulting from naval sonar use, live firing exercises, and certain offshore construction projects such as wind farms, may be harmful to dolphins, increasing stress, damaging hearing, and causing decompression sickness by forcing them to surface too quickly to escape the noise.
In captivity
The renewed popularity of dolphins in the 1960s resulted in the appearance of many dolphinaria around the world, making dolphins accessible to the public. Criticism and animal welfare laws forced many to close, although hundreds still exist around the world. In the United States, the best known are the SeaWorld marine mammal parks.
In the Middle East the best known are Dolphin Bay at Atlantis, The Palm and the Dubai Dolphinarium.
Various species of dolphins are kept in captivity. These small cetaceans are more often than not kept in theme parks, such as SeaWorld, commonly known as a dolphinarium. Bottlenose dolphins are the most common species of dolphin kept in dolphinariums as they are relatively easy to train, have a long lifespan in captivity and have a friendly appearance. Hundreds if not thousands of bottlenose dolphins live in captivity across the world, though exact numbers are hard to determine. Other species kept in captivity are spotted dolphins, false killer whales, and common dolphins, Commerson's dolphins, as well as rough-toothed dolphins, but all in much lower numbers than the bottlenose dolphin. There are also fewer than ten pilot whales, Amazon river dolphins, Risso's dolphins, spinner dolphins, or tucuxi in captivity. Two unusual and very rare hybrid dolphins, known as wholphins, are kept at the Sea Life Park in Hawaii, which is a cross between a bottlenose dolphin and a false killer whale. Also, two common/bottlenose hybrids reside in captivity: one at Discovery Cove and the other at SeaWorld San Diego.
Orca are well known for their performances in shows, but the number of orcas kept in captivity is very small, especially when compared to the number of bottlenose dolphins, with only 44 captive orca being held in aquaria as of 2012. The orca's intelligence, trainability, striking appearance, playfulness in captivity and sheer size have made it a popular exhibit at aquaria and aquatic theme parks. From 1976 to 1997, 55 whales were taken from the wild in Iceland, 19 in Japan and three in Argentina. These figures exclude animals that died during capture. Live captures fell dramatically in the 1990s, and by 1999, about 40% of the 48 animals on display in the world were captive-born.
In history and mythology
In Greek myths, they were seen invariably as helpers of humankind. Dolphins also seem to have been important to the Minoans, judging by artistic evidence from the ruined palace at Knossos. Dolphins are common in Greek mythology, and many coins from ancient Greece have been found which feature a man, a boy or a deity riding on the back of a dolphin. The Ancient Greeks welcomed dolphins; spotting dolphins riding in a ship's wake was considered a good omen. In both ancient and later art, Cupid is often shown riding a dolphin. A dolphin rescued the poet Arion from drowning and carried him safe to land, at Cape Matapan, a promontory forming the southernmost point of the Peloponnesus. There was a temple to Poseidon and a statue of Arion riding the dolphin.
Dionysus was once captured by Etruscan pirates who mistook him for a wealthy prince they could ransom. After the ship set sail Dionysus invoked his divine powers, causing vines to overgrow the ship where the mast and sails had been. He turned the oars into serpents, so terrifying the sailors that they jumped overboard, but Dionysus took pity on them and transformed them into dolphins so that they would spend their lives providing help for those in need. Dolphins were also the messengers of Poseidon and sometimes did errands for him as well. Dolphins were sacred to both Aphrodite and Apollo.
Dolphins are sometimes used as symbols, for instance in heraldry. When heraldry developed in the Middle Ages, not much was known about the biology of the dolphin and it was often depicted as a sort of fish. Traditionally, the dolphins in heraldry still may take after this notion, sometimes showing the dolphin skin covered with fish scales.
In the Middle Ages, the dolphin became an important heraldic element in the coats of arms of several European noble families, the most noticeable being those of the Dauphin de Viennois (later Dauphin of France) through which it passed to the Counts of Forez, Albon and other French families, as well as several branches of the Bourbon family (Count of Montpensier, Count of Beaujolais, among others) the Pandolfini of Florence, and the Delfini of Venice and Rome also used the dolphin as their "canting" armories. In the 19th century, Joseph Bonaparte adopted a dolphin in his coat of arms as King of Naples and Sicily.
In contemporary days, a dolphin is still used in the coat of arms of many cities, as well as in the coat of arms of Anguilla and the coat of arms of Romania, and the coat of arms of Barbados has a dolphin supporter. Cardinal Angelo Amato, prefect of the Congregation of the Causes of Saints, a dolphin in his coat of arms, as well as Cardinal Godfried Danneels, former Metropolitan Archbishop of Mechelen-Brussels.
| Biology and health sciences | Toothed whale | Animals |
305715 | https://en.wikipedia.org/wiki/Lupinus | Lupinus | Lupinus, commonly known as lupin, lupine, or regionally bluebonnet, is a genus of plants in the legume family Fabaceae. The genus includes over 199 species, with centres of diversity in North and South America. Smaller centres occur in North Africa and the Mediterranean. They are widely cultivated, both as a food source and as ornamental plants, but are invasive to some areas.
Description
The species are mostly herbaceous perennial plants tall, but some are annual plants and a few are shrubs up to tall. An exception is the chamis de monte (Lupinus jaimehintonianus) of Oaxaca in Mexico, which is a tree up to tall.
Lupins have soft green to grey-green leaves which may be coated in silvery hairs, often densely so. The leaf blades are usually palmately divided into five to 28 leaflets, or reduced to a single leaflet in a few species of the southeastern United States and eastern South America.
The flowers are produced in dense or open whorls on an erect spike, each flower long. The pea-like flowers have an upper standard, or banner, two lateral wings, and two lower petals fused into a keel. The flower shape has inspired common names such as bluebonnets and quaker bonnets.
The fruit is a pod containing several seeds. The seeds contain alkaloids which lend them a bitter taste.
Taxonomy
The genus Lupinus L. and, in particular, its North American species were divided by Sereno Watson (1873) into three sections: Lupinus, Platycarpos, and Lupinnelus. Differences in habitat and in the number of ovules were the basis for this classification. A majority of the perennial and annual species from the American continent described by Watson were referred to Lupinus. Some annual species with two ovules in the ovary and two seeds in the pod (L. densiflorus, L. microcarpus, etc.) were attributed to the Platycarpos section. Section Lupinnelus consisted of one species (L. uncialis), with axillary and solitary flowers, scarcely reflexed banner, and also with two ovules in the ovary.
While Watson's work was predominantly based on study of North American species, the later research of Ascherson and Graebner (1907) extended his principle of classification to cover all lupins from the Eastern and Western Hemispheres, also using number of ovules (seedbuds) in the ovary (and thus of seeds in the pod) as the criterion for this division. They described two subgenera, Eulupinus and Platycarpos. Most of the described species were referred to subgen. A. Eulupinus. Subgen. B. Platycarpos included several annual species from the Eastern Hemisphere with two seedbuds and seeds in the bean (the same species, as the one specified by S. Watson).
A current schema retains this distinction, but uses the nomenclature for the subgenera of Platycarpos and Lupinus. In this schema, subgenus Platycarpos (S.Wats.) Kurl. contains perennial and annual species from the Western Hemisphere, with a minimum two or more ovules or seedbuds. Subgenus Lupinus consists of 12 species from Africa and the Mediterranean, with a minimum of four ovules or seedbuds.
The taxonomy of Lupinus has always been confusing. How many distinct species exist or how they might be organized within the genus is not clear. The plants are variable and the taxa are not always distinct from one another. Some American taxa have been described as complexes rather than separate species. Estimates of the number of lupine species generally fall between 200 and 500. One authority places the estimate at approximately 267 species worldwide. Currently, two subgenera are recognized.
Subgenus Platycarpos
The ovary contains two and more ovules or seedbuds. The seed are predominantly small-sized, with an underdeveloped embryo and small amount of endosperm. Cotyledons are small-sized, with long caulicles. The first pair of true leaves is alternate. The stem is predominantly naked with waxen coating. Dominating is the monopodial type of branching. Leaflets are smooth, with waxen coating or slight pubescence, predominantly narrow. Pods are flat or orbicular, with two or more seeds. Represented by frutcuilose, fruticose and herbaceous perennial forms, or less often annual ones. Plants are cross-pollinated. Chromosome number 2n is either 36, 48, or 96. This subgenus is distributed throughout North, Central and South America, predominantly in the mining systems of the Andes and Cordillera. Some species are cultivated (L. mutabilis, L. polyphyllus). This subgenus includes several hundred species, requiring further analysis of their authenticity.
It comprises the following species:
Lupinus aberrans C.P. Sm.
Lupinus abramsii C.P. Sm. – Abrams' lupine
Lupinus acopalcus C.P. Sm.
Lupinus adinoanthus C.P. Sm.
Lupinus adsurgens Drew – Drew's silky lupine
Lupinus affinis J. Agardh – fleshy lupine
Lupinus agardhianus A. Heller
Lupinus alaristatus C.P. Sm.
Lupinus albert-smithianus C.P. Sm.
Lupinus albescens Hook. & Arn. – hoary lupine
Lupinus albicaulis Douglas – sickle-keel lupine
Lupinus albifrons Benth. – silver bush lupine
var. albifrons Benth.
var. douglasii (J. Agardh) C. P. Sm.
var. hallii (Abrams) Isely
Lupinus albopilosus A. Heller
Lupinus albosericeus C.P. Sm.
Lupinus alcis-montis C.P. Sm.
Lupinus aliamandus C.P. Sm.
Lupinus aliattenuatus C.P. Sm.
Lupinus alibicolor C.P. Sm.
Lupinus aliceae C.P. Sm.
Lupinus alilatissimus C.P. Sm.
Lupinus alinanus C.P. Sm.
Lupinus alipatulus C.P. Sm.
Lupinus alirevolutus C.P. Sm.
Lupinus alivillosus C.P. Sm.
Lupinus allargyreius C.P. Sm.
Lupinus alopecuroides Desr.
Lupinus alpestris A. Nelson
Lupinus altimontanus C.P. Sm.
Lupinus altiplani C.P. Sm.
Lupinus amabayensis C.P. Sm.
Lupinus amandus C.P. Sm.
Lupinus amboensis C.P. Sm.
Lupinus ammophilus Greene
var. ammophilus Greene
var. crassus (Payson) Isely
Lupinus amnis-otuni C.P. Sm.
Lupinus ampaiensis C.P. Sm.
Lupinus amphibius Suksd.
Lupinus ananeanus Ulbr.
Lupinus anatolicus W. Święcicki & W. K. Święcicki
Lupinus andersonii S. Watson – Anderson's lupine
Lupinus andicola Gillies
Lupinus andinus Rose ex J. F. Macbr.
Lupinus angustiflorus Eastw. – narrowflower lupine
Lupinus antensis C.P. Sm.
Lupinus antiplani C. P. Sm.
Lupinus antoninus Eastw. – Anthony Peak lupine
Lupinus apertus A. Heller
Lupinus appositus C.P. Sm.
Lupinus arboreus Sims – yellow bush lupin, tree lupine
Lupinus arbustus Lindl. – longspur lupine
subsp. arbustus Lindl.
subsp. neolaxiflorus D.B.Dunn
subsp. pseudoparviflorus (Rydb.) D.B.Dunn
Lupinus arbutosocius C.P. Sm.
Lupinus archeranus C.P. Sm.
Lupinus arcticus S. Watson – Arctic lupine
subsp. arcticus S. Watson
subsp. subalpinus (Piper & Robinson)D.B.Dunn
Lupinus arenarius Gardner
Lupinus arequipensis C.P. Sm.
Lupinus argenteus Pursh – silvery lupine
var. argentatus (Rydb.) Barneby
var. argenteus Pursh
var. argophyllus (A. Gray) S. Watson
var. depressus (Rydb.) C. L. Hitchc.
var. fulvomaculatus (Payson) Barneby
var. heteranthus (S. Watson) Barneby – Kellogg's spurred lupine
var. hillii (Greene) Barneby
var. holosericeus (Torr. & A.Gray) Barneby
var. montigenus (A. Heller) Barneby
var. palmeri (S.Watson) Barneby
var. rubricaulis (Greene) S. L. Welsh
var. utahensis (S.Watson) Barneby
Lupinus argurocalyx C.P. Sm.
Lupinus aridorum McFarlin ex Beckner – scrub lupine
Lupinus aridulus C.P. Sm.
Lupinus aridus Lindl.
Lupinus ariste-josephii C.P. Sm.
Lupinus arizelus C.P. Sm.
Lupinus arizonicus (S. Watson) S. Watson
subsp. arizonicus (S. Watson) S. Watson – Arizona lupine
subsp. sonorensis Christian & D. Dunn – Sonora lupine
Lupinus arvensi-plasketti C.P. Sm.
Lupinus arvensis Benth.
Lupinus asa-grayanus C.P. Sm.
Lupinus aschenbornii S. Schauer
Lupinus asplundianus C.P. Sm.
Lupinus asymbepus C.P. Sm.
Lupinus atropurpureus C.P. Sm.
Lupinus attenuatus Gardner
Lupinus aureonitens Hook. & Arn.
Lupinus austrobicolor C.P. Sm.
Lupinus austrohumifusus C.P. Sm.
Lupinus austrorientalis C.P. Sm.
Lupinus austrosericeus C.P. Sm.
Lupinus ballianus C.P. Sm.
Lupinus bandelierae C.P. Sm.
Lupinus bangii Rusby
Lupinus barbatilabius C.P. Sm.
Lupinus barkeri Lindl.
Lupinus bartlettianus C.P. Sm.
Lupinus benthamii A. Heller
Lupinus bi-inclinatus C.P. Sm.
Lupinus bicolor Lindl. – miniature lupine, bicolor lupine, Lindley's annual lupine
subsp. bicolor Lindl.
subsp. microphyllus (S. Watson) D. B. Dunn
subsp. pipersmithii (A. Heller) D. B. Dunn
subsp. umbellatus (Greene) D. B. Dunn
Lupinus bingenensis Suksd. – Bingen lupine
Lupinus blaisdellii Eastw.
Lupinus bogotensis Benth.
Lupinus bolivianus C.P. Sm.
Lupinus bombycinocarpus C.P. Sm.
Lupinus bonplandius C.P. Sm.
Lupinus boyacensis C.P. Sm.
Lupinus brachypremnon C.P. Sm.
Lupinus bracteolaris Desr.
Lupinus brandegeei Eastw.
Lupinus brevecuneus C.P. Sm.
Lupinus brevicaulis S. Watson – shortstem lupine
Lupinus brevior (Jeps.) Christian & D.B. Dunn
Lupinus breviscapus Ulbr.
Lupinus breweri A. Gray – Brewer's lupine
Lupinus bryoides C.P. Sm.
Lupinus buchtienii Rusby
Lupinus burkartianus C.P. Sm.
Lupinus burkei S. Watson – Burke's lupine
Lupinus burkeri Lindl.
Lupinus caballoanus B.L. Turner
Lupinus cachupatensis C.P. Sm.
Lupinus cacuminis Standl.
Lupinus caeruleus A. Heller
Lupinus caesius Eastw.
Lupinus caespitosus Torr. & A. Gray – stemless dwarf lupine
Lupinus calcensis C.P. Sm.
Lupinus caldasensis C.P. Sm.
Lupinus camiloanus C.P. Sm.
Lupinus campestris Schltdl. & Cham.
Lupinus carazensis Ulbr.
Lupinus carchiensis C.P. Sm.
Lupinus cardenasianus C.P. Sm.
Lupinus carhuamayus C.P. Sm.
Lupinus carlos-ochoae C.P. Sm.
Lupinus carpapaticus C.P. Sm.
Lupinus carrikeri C.P. Sm.
Lupinus caucensis C.P. Sm.
Lupinus cavicaulis C.P. Sm.
Lupinus ccorilazensis Vargas ex C. P. Smith
Lupinus celsimontanus C.P. Sm.
Lupinus cervinus Kellogg – Santa Lucia lupine
Lupinus cesar-vargasii C.P. Sm.
Lupinus cesaranus C.P. Sm.
Lupinus chachas C.P. Sm.
Lupinus chamissonis Eschsch. – Chamisso bush lupine
Lupinus chavanillensis (J.F. Macbr.) C.P. Sm.
Lupinus chipaquensis C.P. Sm.
Lupinus chlorolepis C.P. Sm.
Lupinus chocontensis C.P. Sm.
Lupinus chongos-bajous C.P. Sm.
Lupinus christinae A. Heller
Lupinus chrysanthus Ulbr.
Lupinus chrysocalyx C.P. Sm.
Lupinus chumbivilcensis C.P. Sm.
Lupinus citrinus Kellogg – orange lupine
Lupinus clarkei Oerst.
Lupinus cochapatensis C.P. Sm.
Lupinus colcabambensis C.P. Sm.
Lupinus collinus (Greene) A. Heller
Lupinus colombiensis C.P. Sm.
Lupinus compactiflorus Rose
Lupinus comptus Benth.
Lupinus concinnus J. Agardh
subsp. concinnus J. Agardh
subsp. orcuttii (S.Watson) D.B.Dunn
Lupinus condensiflorus C.P. Sm.
Lupinus confertus Kellogg
Lupinus congdonii (C.P. Sm.) D.B. Dunn
Lupinus conicus C.P. Sm.
Lupinus constancei T.W. Nelson & J.P. Nelson – Lassics lupine
Lupinus convencionensis C.P. Sm.
Lupinus cookianus C.P. Sm.
Lupinus coriaceus Benth.
Lupinus costaricensis D.B. Dunn
Lupinus cotopaxiensis C.P. Sm.
Lupinus couthouyanus C.P. Sm.
Lupinus covillei Greene – shaggy lupine
Lupinus crassulus Greene
Lupinus crassus Payson
Lupinus croceus Eastw. – saffron-flowered lupine
Lupinus crotalarioides Benth.
Lupinus crucis-viridis C.P. Sm.
Lupinus cuatrecasasii C.P. Sm.
Lupinus culbertsonii Greene
subsp. culbertsonii Greene
subsp. hypolasius (Greene) B.J.Cox
Lupinus cumulicola Small
Lupinus cusickii S. Watson
subsp. abortivus (Greene) B.J.Cox
subsp. brachypodus (Piper) B.J.Cox
subsp. cusickii S. Watson
Lupinus cuspidatus Rusby
Lupinus cuzcensis C.P. Sm.
Lupinus cymboides C.P. Sm.
Lupinus czermakii Briq. & Hochr.
Lupinus dalesiae Eastw. – Quincy lupine
Lupinus decemplex C.P. Sm.
Lupinus decurrens Gardner
Lupinus deflexus Congdon
Lupinus delicatulus Sprague & Riley
Lupinus densiflorus Benth. – dense-flowered lupin
subsp. densiflorus Benth.
subsp. lacteus (Kellogg) R.M.Beauch.
Lupinus depressus Rydb.
Lupinus diasemus C.P. Sm.
Lupinus diehlii M.E. Jones
Lupinus diffusus Nutt. – spreading lupine, Oak Ridge lupine, sky-blue lupine
Lupinus disjunctus C.P. Sm.
Lupinus diversalpicola C.P. Sm.
Lupinus dorae C.P. Sm.
Lupinus dotatus C.P. Sm.
Lupinus duranii Eastw. – Mono Lake lupine
Lupinus dusenianus C.P. Sm.
Lupinus eanophyllus C.P. Sm.
Lupinus edysomatus C.P. Sm.
Lupinus egens C.P. Sm.
Lupinus elaphoglossum Barneby
Lupinus elatus I.M. Johnst. – tall silky lupine
Lupinus elegans Kunth – elegant lupine
Lupinus elegantulus Eastw.
Lupinus ellsworthianus C.P. Sm.
Lupinus elmeri Greene – Elmer's lupine
Lupinus eramosus C.P. Sm.
Lupinus erectifolius C.P. Sm.
Lupinus eremonomus C.P. Sm.
Lupinus eriocalyx (C.P. Sm.) C.P. Sm.
Lupinus eriocladus Ulbr.
Lupinus evermannii Rydb.
Lupinus espinarensis C.P. Sm.
Lupinus exaltatus Zucc.
Lupinus excubitus M.E. Jones – grape soda lupine
subsp. austromontanus (A.Heller) R.M.Beauch.
subsp. excubitus M.E. Jones
Lupinus exochus C.P. Sm.
Lupinus expetendus C.P. Sm.
Lupinus extrarius C.P. Sm.
Lupinus falsomutabilis C.P. Sm.
Lupinus falsoprostratus C.P. Sm.
Lupinus falsorevolutus C.P. Sm.
Lupinus famelicus C.P. Sm.
Lupinus fiebrigianus Ulbr.
Lupinus fieldii J.F. Macbr.
Lupinus fissicalyx A. Heller
Lupinus flavoculatus A. Heller
Lupinus foliolosus Benth.
Lupinus formosus Greene – summer lupine
var. bridgesii (S.Watson) Greene
var. formosus Greene
Lupinus fragrans A. Heller
Lupinus francis-whittieri C.P. Sm.
Lupinus fratrum C.P. Sm.
Lupinus fulcratus Greene
Lupinus gachetensis C.P. Sm.
Lupinus garfieldensis C.P. Sm.
Lupinus gaudichaudianus C.P. Sm.
Lupinus gayanus C.P. Sm.
Lupinus gentryanus C.P. Sm.
Lupinus geophilus Rose
Lupinus gibertianus C.P. Sm.
Lupinus giganteus Rose
Lupinus glabratus J. Agardh
Lupinus goodspeedii J.F. Macbr.
Lupinus gormanii Piper
Lupinus gracilentus Greene
Lupinus grayi S. Watson – Sierra lupine
Lupinus grauensis C.P. Sm.
Lupinus grisebachianus C.P. Sm.
Lupinus guadalupensis C.P. Sm. – Guadalupe Island lupine
Lupinus guaraniticus (Hassl.) C.P. Sm.
Lupinus guascensis C.P. Sm.
Lupinus guggenheimianus Rusby
Lupinus hamaticalyx C.P. Sm.
Lupinus hartmannii C.P. Sm.
Lupinus hartwegii Lindl.
Lupinus haughtianus C.P. Sm.
Lupinus hautcarazensis C.P. Sm.
Lupinus havardii S. Watson
Lupinus hazenanus C.P. Sm.
Lupinus hendersonii Eastw.
Lupinus heptaphyllus (Vell.) Hassl.
Lupinus herreranus C.P. Sm.
Lupinus herzogii Ulbr.
Lupinus hieronymii C.P. Sm.
Lupinus hilarianus Benth.
Lupinus hillii Greene
Lupinus hinkleyorum C.P. Sm.
Lupinus hintoniorum B.L. Turner
Lupinus hirsutissimus Benth. – stinging lupine
Lupinus holmgrenianus C.P. Sm. – Holmgren's lupine
Lupinus honoratus C.P. Sm.
Lupinus horizontalis A. Heller
Lupinus hornemanni J. Agardh
Lupinus hortonianus C.P. Sm.
Lupinus hortorum C.P. Sm.
Lupinus howard-scottii C.P. Sm.
Lupinus howardii M.E. Jones
Lupinus huachucanus M.E. Jones
Lupinus huancayoensis C.P. Sm.
Lupinus huariacus C.P. Sm.
Lupinus huaronensis J.F. Macbr.
Lupinus huigrensis Rose ex C. P. Sm.
Lupinus humifusus Sessé & Moc. ex G. Don
Lupinus hyacinthinus C.F. Baker – San Jacinto lupine
Lupinus hybridus Lem.
Lupinus ignobilis C.P. Sm.
Lupinus imminutus C.P. Sm.
Lupinus indigoticus Eastw.
Lupinus inflatus C.P. Sm.
Lupinus insignis C.P. Sm.
Lupinus insulae C.P. Sm.
Lupinus interruptus Benth.
Lupinus intortus C.P. Sm.
Lupinus inusitatus C.P. Sm.
Lupinus involutus C.P. Sm.
Lupinus inyoensis A. Heller
Lupinus isabelianus Eastw.
Lupinus jahnii Rose ex Pittier
Lupinus jaimehintoniana B.L. Turner
Lupinus james-westii C.P. Sm.
Lupinus jean-julesii C.P. Sm.
Lupinus jelskianus C.P. Sm.
Lupinus johannis-howellii C.P. Sm.
Lupinus jonesii Rydb.
Lupinus jujuyensis C.P. Sm.
Lupinus juninensis C.P. Sm.
Lupinus kalenbornorum C.P. Sm.
Lupinus kellermanianus C.P. Sm.
Lupinus kerrii Eastw.
Lupinus killipianus C.P. Sm.
Lupinus kingii S. Watson
Lupinus klamathensis Eastw.
Lupinus kunthii J. Agardh
Lupinus kuschei Eastw. – Yukon lupine
Lupinus lacus C.P. Sm.
Lupinus laetus Wooton & Standl.
Lupinus laevigatus Benth.
Lupinus lagunae-negrae C.P. Sm.
Lupinus lanatocarpus C.P. Sm.
Lupinus lanatus Benth.
Lupinus lapidicola A. Heller – Mt. Eddy lupine
Lupinus latifolius J. Agardh
subsp. dudleyi (C.P.Sm.) P.Kenney & D.B.Dunn
subsp. latifolius J. Agardh
var. latifolius J. Agardh – broadleaf lupine
var. barbatus – Klamath lupine, bearded lupine
subsp. leucanthus (Rydb.)P.Kenney & D.B.Dunn
subsp. longipes (Greene) P.Kenney & D.B.Dunn
subsp. parishii (C.P.Sm.) P.Kenney & D.B.Dunn
subsp. viridifolius (A.Heller) P.Kenney & D.B.Dunn
Lupinus laudandrus C.P. Sm.
Lupinus lechlerianus C.P. Sm.
Lupinus ledigianus C.P. Sm.
Lupinus lelandsmithii Eastw.
Lupinus lemmonii C.P. Sm.
Lupinus lepidus Lindl. – prairie lupine
var. aridus (Douglas) Jeps.
var. confertus (Kellogg) C. P. Sm.
var. lepidus Lindl.
var. lobbii (A. Gray ex S. Watson) C. L. Hitchc.
var. sellulus (Kellogg) Barneby
var. utahensis (S. Watson) C. L. Hitchc.
Lupinus leptocarpus Benth.
Lupinus leptophyllus Cham. & Schltdl.
Lupinus lespedezoides C.P. Sm.
Lupinus leucophyllus Lindl. – woolly-leaf lupine
Lupinus lilacinus A. Heller
Lupinus lindenianus C.P. Sm.
Lupinus lindleyanus J. Agardh
Lupinus linearis Desr.
Lupinus littoralis Lindl. – seashore lupine
Lupinus lobbianus C.P. Sm.
Lupinus longifolius (S. Watson) Abrams – longleaf bush lupine
Lupinus lorenzensis C.P. Sm.
Lupinus ludovicianus Greene – San Luis Obispo County Lupine
Lupinus luetzelburgianus C.P. Sm.
Lupinus luteolus Kellogg – butter lupine, pale yellow lupine
Lupinus lutescens C.P. Sm.
Lupinus lutosus A. Heller
Lupinus lyallii A. Gray
subsp. alcis-temporis (C.P. Sm.) B.J.Cox
subsp. lyallii A. Gray – Lyall's lupine
subsp. minutifolius (Eastw.) B.J.Cox
subsp. washoensis (A.Heller) B.J.Cox
Lupinus macbrideanus C.P. Sm.
Lupinus macranthus Rose
Lupinus maculatus Rydb.
Lupinus madrensis Seem.
Lupinus magdalenensis C.P. Sm.
Lupinus magnificus M.E. Jones
Lupinus magniflorus C.P. Sm.
Lupinus magnistipulatus Planchuelo & D.B. Dunn
Lupinus malacophyllus Greene
Lupinus malacotrichus C.P. Sm.
Lupinus maleopinatus C.P. Sm.
Lupinus mandonanus C.P. Sm.
Lupinus mantaroensis C.P. Sm.
Lupinus marinensis Eastw.
Lupinus mariposanus Eastw.
Lupinus martensis C.P. Sm.
Lupinus martinetianus (C.P. Sm.) C.P. Sm.
Lupinus mathewsianus C.P. Sm.
Lupinus matucanicus Ulbr.
Lupinus meionanthus A. Gray
Lupinus melaphyllus C.P. Sm.
Lupinus menziesii J. Agardh
Lupinus meridanus C.P. Sm.
Lupinus metensis C.P. Sm.
Lupinus mexicanus Lag.
Lupinus michelianus C. P. Sm.
Lupinus microcarpus Sims
var. densiflorus
var. microcarpus – wide-bannered lupin, chick lupin
Lupinus microphyllus Desr.
Lupinus minimus Hook.
Lupinus mirabilis C.P. Sm.
Lupinus misticola Ulbr.
Lupinus mollendoensis Ulbr.
Lupinus mollis A. Heller
Lupinus monensis Eastw.
Lupinus monserratensis C.P. Sm.
Lupinus montanus Kunth
subsp. glabrior (S.Watson) D.B.Dunn & Harmon
subsp. montanus Kunth
subsp. montesii (C.P.Sm.) D.B.Dunn & Harmon
Lupinus monticola Rydb.
Lupinus montigenus A. Heller
Lupinus moritzianus Kunth
Lupinus mucronulatus Howell
Lupinus muelleri Standl.
Lupinus multiflorus Desr.
Lupinus munzianus C.P. Sm.
Lupinus munzii Eastw.
Lupinus mutabilis Sweet – Andean lupin, pearl lupin, South American lupin, tarwi, tarhui, chocho
Lupinus nanus Benth. – dwarf lupin, field lupin, sky lupin, Douglas' annual lupin
Lupinus navicularius A. Heller
Lupinus nehmadae C.P. Sm.
Lupinus neocotus C.P. Sm.
Lupinus neomexicanus Greene
Lupinus nepubescens C.P. Sm.
Lupinus nevadensis A. Heller – Nevada lupine
Lupinus niederleinianus C.P. Sm.
Lupinus nipomensis Eastw. – Nipomo Mesa lupine
Lupinus niveus S. Watson
Lupinus nonoensis C.P. Sm.
Lupinus nootkatensis Sims – Nootka lupin
Lupinus notabilis C.P. Sm.
Lupinus nubigenus Kunth
Lupinus nubilorum C.P. Sm.
Lupinus obscurus C.P. Sm.
Lupinus obtusilobus A. Heller – bluntlobe lupine
Lupinus ochoanus C.P. Sm.
Lupinus ochroleucus Eastw.
Lupinus odoratus A. Heller – royal Mojave lupin
Lupinus onustus S. Watson – Plumas lupine
Lupinus opertospicus C.P. Sm.
Lupinus oquendoanus C.P. Sm.
Lupinus oreganus A. Heller – Oregon lupin
Lupinus oreophilus Phil.
Lupinus ornatus Lindl.
Lupinus oscar-haughtii C.P. Sm.
Lupinus ostiofluminis C.P. Sm.
Lupinus otto-buchtienii C.P. Sm.
Lupinus otto-kuntzeanus C.P. Sm.
Lupinus otuzcoensis C.P. Sm.
Lupinus ovalifolius Benth.
Lupinus pachanoanus C.P. Sm.
Lupinus pachitensis C.P. Sm.
Lupinus pachylobus Greene
Lupinus padre-crowleyi C.P. Sm. – DeDecker's lupine, Father Crowley's lupine
Lupinus pallidus Brandegee
Lupinus paniculatus Desr.
Lupinus paraguariensis Chodat & Hassl.
Lupinus paranensis C.P. Sm.
Lupinus paruroensis C.P. Sm.
Lupinus parviflorus Hook. & Arn. – lodgepole lupin
subsp. myrianthus (Greene) Harmon
subsp. parviflorus Hook. & Arn.
Lupinus parvifolius Gardner
Lupinus pasachoensis C.P. Sm.
Lupinus pasadenensis Eastw.
Lupinus patulus C.P. Sm.
Lupinus paucartambensis C.P. Sm.
Lupinus paucovillosus C.P. Sm.
Lupinus paynei Davidson
Lupinus pearceanus C.P. Sm.
Lupinus pendentiflorus C.P. Sm.
Lupinus peirsonii H. Mason – Peirson's lupine, long lupine
Lupinus penlandianus C.P. Sm.
Lupinus perblandus C.P. Sm.
Lupinus perbonus C.P. Sm.
Lupinus perennis L. – wild perennial lupine, sundial lupine, Indian beet, old maid's bonnets
subsp. gracilis (Nutt.) D.B.Dunn
subsp. occidentalis S. Watson
subsp. perennis L.
Lupinus perglaber Eastw.
Lupinus perissophytus C.P. Sm.
Lupinus persistens Rose
Lupinus peruvianus Ulbr.
Lupinus philippianus C.P. Sm.
Lupinus physodes Douglas
Lupinus pickeringii A. Gray
Lupinus pilosellus Eastw.
Lupinus pilosissimus M. Martens & Galeotti
Lupinus pinguis Ulbr.
Lupinus pipersmithianus J.F. Macbr.
Lupinus pisacensis C.P. Sm.
Lupinus piurensis C.P. Sm.
Lupinus platamodes C.P. Sm.
Lupinus plattensis S. Watson
Lupinus platyptenus C.P. Sm.
Lupinus polycarpus Greene – smallflower lupin
Lupinus polyphyllus Lindl. – largeleaf lupine, bigleaf lupine, garden lupin, many-leaved lupine
var. burkei (S. Watson) C. L. Hitchc.
var. humicola (A.Nelson) Barneby
var. pallidipes (A. Heller) C. P. Sm.
var. polyphyllus Lindl.
var. prunophilus (M. E. Jones) L. Ll. Phillips
Lupinus poopoensis C.P. Sm.
Lupinus popayanensis C.P. Sm.
Lupinus potosinus Rose
Lupinus praealtus C.P. Sm.
Lupinus praestabilis C.P. Sm.
Lupinus praetermissus C.P. Sm.
Lupinus pratensis A.Heller – Inyo Meadow lupine
Lupinus pringlei Rose
Lupinus proculaustrinus C.P. Sm.
Lupinus prostratus J. Agardh
Lupinus protrusus C.P. Sm.
Lupinus prouvensalanus C.P. Sm.
Lupinus prunophilus M.E. Jones – hairy bigleaf lupin
Lupinus pseudopolyphyllus C.P. Sm.
Lupinus pseudotsugoides C.P. Sm.
Lupinus pubescens Benth.
Lupinus pucapucensis C.P. Sm.
Lupinus pulloviridus C.P. Sm.
Lupinus pulvinaris Ulbr.
Lupinus punto-reyesensis C.P. Sm.
Lupinus puracensis C.P. Sm.
Lupinus purdieanus C.P. Sm.
Lupinus pureriae C.P. Sm.
Lupinus purosericeus C.P. Sm.
Lupinus pusillus Pursh – rusty lupine or dwarf lupine
subsp. intermontanus (A.Heller) D.B.Dunn
subsp. pusillus Pursh
Lupinus puyupatensis C.P. Sm.
Lupinus pycnostachys C.P. Sm.
Lupinus quellomayus C.P. Sm.
Lupinus quitensis C.P. Sm.
Lupinus radiatus C.P. Sm.
Lupinus ramosissimus Benth.
Lupinus reflexus Rose
Lupinus regalis Bergmans
Lupinus regnellianus C.P. Sm.
Lupinus reineckianus C.P. Sm.
Lupinus reitzii Burkart ex M. Pinheiro & Miotto
Lupinus retrorsus L.F. Hend.
Lupinus revolutus C.P. Sm.
Lupinus richardianus C.P. Sm.
Lupinus rimae Eastw.
Lupinus rivularis Lindl. – riverbank lupin
Lupinus romasanus Ulbr.
Lupinus roseolus Rydb.
Lupinus roseorum C.P. Sm.
Lupinus rotundiflorus M.E. Jones
Lupinus rowleeanus C.P. Sm.
Lupinus ruber A. Heller
Lupinus rubriflorus Planchuelo
Lupinus ruizensis C.P. Sm.
Lupinus rupestris Kunth
Lupinus rusbyanus C.P. Sm.
Lupinus russellianus C.P. Sm.
Lupinus sabinianus Lindl.
Lupinus sabinii Hook.
Lupinus sabulosus A. Heller
Lupinus salticola Eastw.
Lupinus sandiensis C.P. Sm.
Lupinus santanderensis C.P. Sm.
Lupinus sarmentosus Desr.
Lupinus saxatilis Ulbr.
Lupinus saxosus Howell – rock lupine
Lupinus schwackeanus C.P. Sm.
Lupinus seifrizianus (C.P. Sm.) C.P. Sm.
Lupinus sellowianus Harms
Lupinus sellulus Kellogg
var. lobbii (S.Watson) B.J.Cox
var. sellulus Kellogg
var. ursinus (Eastw.) B.J.Cox
Lupinus semiprostratus C.P. Sm.
Lupinus semperflorens Benth.
Lupinus sericatus Kellogg – Cobb Mountain lupine
Lupinus sericeus Pursh – Pursh's silky lupin
var. barbiger (S.Watson) S.L.Welsh
var. sericeus Pursh
Lupinus setifolius Planchuelo & D.B. Dunn
Lupinus shastensis Lupinus albicaulis
Lupinus shockleyi S. Watson – purple desert lupine
Lupinus sierrae-blancae Wooton & Standl.
subsp. aquilinus (Wooton & Standl.) L.S.Fleak & D.B.Dunn
subsp. sierrae-blancae Wooton & Standl.
Lupinus simonsianus C.P. Sm.
Lupinus simulans Rose
Lupinus sinaloensis C.P. Sm.
Lupinus sitgreavesii S. Watson
Lupinus smithianus Kunth
Lupinus solanagrorum C.P. Sm.
Lupinus sonomensis A. Heller
Lupinus soratensis Rusby
Lupinus soukupianus C. P. Smith ex J. F. Macbr.
Lupinus sparsiflorus Benth. – desert lupin, Coulter's lupin, Mojave lupin
Lupinus spectabilis Hoover – shaggyhair lupine
Lupinus splendens Rose
Lupinus spragueanus C.P. Sm.
Lupinus staffordiae C.P. Sm.
Lupinus stipulatus J. Agardh
Lupinus stiversii Kellogg – harlequin annual lupine
Lupinus storkianus C.P. Sm.
Lupinus subacaulis Griseb.
Lupinus subcarnosus Hook. – buffalo clover
Lupinus subcuneatus C.P. Sm.
Lupinus subhamatus C.P. Sm.
Lupinus subinflatus C.P. Sm.
Lupinus sublanatus Eastw.
Lupinus submontanus Rose
Lupinus subsessilis Benth.
Lupinus subtomentosus C.P. Sm.
Lupinus subvexus C.P. Sm.
Lupinus succulentus K. Koch – succulent lupin, arroyo lupin, hollowleaf annual lupin
Lupinus sufferrugineus Rusby
Lupinus suksdorfii Robinson
Lupinus sulphureus Douglas
subsp. kincaidii (Suksd.) L. Ll. Phillips – Kincaid's lupin
subsp. subsaccatus (Suksd.) L. Ll. Phillips
subsp. sulphureus Douglas – sulphur lupin, sulphur-flowered lupin
Lupinus surcoensis C.P. Sm.
Lupinus syriggedes C.P. Sm.
Lupinus tacitus C.P. Sm.
Lupinus tafiensis C.P. Sm.
Lupinus talahuensis C.P. Sm.
Lupinus tamayoanus C.P. Sm.
Lupinus tarapacensis C.P. Sm.
Lupinus tarijensis Ulbr.
Lupinus tarmaensis C.P. Sm.
Lupinus tatei Rusby
Lupinus taurimortuus C.P. Sm.
Lupinus tauris Benth.
Lupinus tayacajensis C.P. Sm.
Lupinus tegeticulatus Eastw.
Lupinus tetracercophorus C.P. Sm.
Lupinus texanus Hook.
Lupinus texensis Hook. – Texas bluebonnet
Lupinus thompsonianus C.P. Sm.
Lupinus tidestromii Greene – Tidestrøm's lupin
var. layneae (Eastw.) Munz
var. tidestromii Greene
Lupinus tolimensis C.P. Sm.
Lupinus tomentosus DC.
Lupinus tominensis Wedd.
Lupinus toratensis C.P. Sm. – warwanzo, lito
Lupinus tracyi Eastw. – Tracy's lupine
Lupinus triananus C.P. Sm.
Lupinus truncatus Hook. & Arn. – collared annual lupine
Lupinus tucumanensis C.P. Sm.
Lupinus ulbrichianus C.P. Sm.
Lupinus uleanus C.P. Sm.
Lupinus ultramontanus C.P. Sm.
Lupinus umidicola C.P. Sm.
Lupinus uncialis S. Watson
Lupinus uncinatus Schltdl.
Lupinus urcoensis C.P. Sm.
Lupinus urubambensis C.P. Sm.
Lupinus valerioi Standl.
Lupinus vallicola A. Heller – open lupin
subsp. apricus (Greene) D.B.Dunn
subsp. vallicola A. Heller
Lupinus vargasianus C.P. Sm.
Lupinus varicaulis C.P. Sm.
Lupinus variicolor Steud. – varied lupin
Lupinus velillensis C.P. Sm.
Lupinus velutinus Benth.
Lupinus venezuelensis C.P. Sm.
Lupinus ventosus C.P. Sm.
Lupinus verbasciformis Sandwith
Lupinus verjonensis C.P. Sm.
Lupinus vernicius Rose
Lupinus viduus C.P. Sm.
Lupinus vilcabambensis C.P. Sm.
Lupinus villosus Willd.
Lupinus visoensis J.F. Macbr.
Lupinus volubilis C.P. Sm.
Lupinus weberbaueri Ulbr.
Lupinus werdermannianus C.P. Sm.
Lupinus westianus Small
var. aridorum (McFarlin ex Beckner) Isely
var. westianus Small
Lupinus whiltoniae Eastw.
Lupinus wilkesianus C.P. Sm.
Lupinus williamlobbii C.P. Sm.
Lupinus williamsianus C.P. Sm.
Lupinus xanthophyllus C.P. Sm.
Lupinus xenophytus C.P. Sm.
Lupinus yanahuancensis C.P. Sm.
Lupinus yarushensis C.P. Sm.
Lupinus ynesiae C.P. Sm.
Subgenus Lupinus
In its current circumscription, subgenus Lupinus includes 12 species from the Mediterranean region and Africa with at least four ovules or seedbuds in the ovary:
Lupinus albus L. 1753 – white lupine
subsp. albus L.
subsp. graecus (Boiss. & Spruner) Franco & P.Silva
subsp. termis (Forsk.) Ponert.
Lupinus angustifolius L. 1753 – blue lupin, narrow-leafed lupin
var. angustifolius L.
var. albopunctatus Kurl. et Stankev.
var. griseomaculatus Kurl. et Stankev.
var. chalybens Kurl. et Stankev.
var. corylinus Kurl. et Stankev.
var. purpureus Kurl. et Stankev.
var. rubidus Kurl. et Stankev.
var. atabekovae Kurl. et Stankev.
var. sparsiusculus Kurl. et Stankev.
var. brunneus Kurl. et Stankev.
var. albosyringeus Taran.
var. albidus Kurl. et Stankev.
var. candidus Kuptzov. et Kurl.
Lupinus atlanticus Gladstones 1974
Lupinus cosentinii Guss. 1828 – sandplain lupin
Lupinus digitatus Forsk. 1775
Lupinus hispanicus Boiss. & Reut. 1842
subsp. bicolor (Merino) Gladst.
subsp. hispanicus Boiss. & Reut.
Lupinus luteus L. 1753 – yellow lupin
var. luteus L.
var. maculosus Kurl. et Stankev.
var. kazimierskii Kurl. et Stankev.
var. arcellus Kurl. et Stankev.
var. sempolovskii (Atab) Kurl. et Stankev.
var. melanospermus Kurl. et Stankev.
var. niger Kurl. et Stankev.
var. cremeus Kurl. et Stankev.
var. leucospermus Kurl. et Stankev.
var. sulphureus (Atab.) Kurl. et Stankev.
var. stepanovae Kurl. et Stankev.
var. ochroleucus Kurl. et Stankev.
var. aurantiacus Kurl. et Stankev.
var. croceus Kurl. et Stankev.
var. aureus Kurl. et Stankev.
var. albicans Kurl. et Stankev.
var. sinskayae Kurl. et Stankev.
Lupinus micranthus Guss. 1828
Lupinus palaestinus Boiss. 1849 – white-grey lupine
Lupinus pilosus Murr. 1774 – blue lupine
Lupinus princei Harms 1901
Lupinus somaliensis Baker f. 1895
Species names with uncertain taxonomic status
The status of the following binomials is unresolved:
Lupinus acaulis Larrañaga
Lupinus achilleaphilus C.P.Sm.
Lupinus acutilobus A.Heller
Lupinus aegr-Aovium C.P.Sm.
Lupinus africanus Lour.
Lupinus agninus Gand.
Lupinus agropyrophilus C.P.Sm.
Lupinus alaimandus C.P.Sm.
Lupinus albicaulis Douglas ex Hook.
Lupinus alicanescens C.P.Sm.
Lupinus aliclementinus C.P.Sm.
Lupinus aliumbellatus C.P.Sm.
Lupinus altissimus Sessé & Moc.
Lupinus alturasensis C.P.Sm.
Lupinus alveorum C.P.Sm.
Lupinus amabilis A.Heller
Lupinus amniculi-cervi C.P.Sm.
Lupinus amniculi-salicis C.P.Sm.
Lupinus amniculi-vulpum C.P.Sm.
Lupinus andersonianus C.P.Sm.
Lupinus anemophilus Greene
Lupinus angustifolius Blanco
Lupinus aphronorus Blank.
Lupinus apodotropis A.Heller
Lupinus aralloius C.P.Sm.
Lupinus arborescens Amabekova & Maisuran
Lupinus arceuthinus Greene
Lupinus argyraeus DC.
Lupinus atacamicus C.P.Sm.
Lupinus aureus J.Agardh
Lupinus axillaris Blank.
Lupinus barkeriae Knowles & Westc.
Lupinus bartolomei M.E.Jones
Lupinus bassett-maguirei C.P.Sm.
Lupinus beaneanus C.P.Sm.
Lupinus biddleii L.F.Hend.
Lupinus bimaculatus Hook. ex D.Don
Lupinus bimaculatus Desr.
Lupinus bivonii C.Presl
Lupinus blankinshipii A.Heller
Lupinus blaschkeanus Fisch. & C.A.Mey.
Lupinus brevior (Jeps.) J.A. Christian & D.B. Dunn
Lupinus brittonii Abrams
Lupinus caespitosus Nutt.
Lupinus californicus K.Koch
Lupinus campbelliae Eastw.
Lupinus campestris Cham. & Schltdl.
Lupinus campestris-florum C.P.Sm.
Lupinus candicans Rydb.
Lupinus canus Hemsl.
Lupinus capitatus Greene
Lupinus capitis-amniculi C.P.Sm.
Lupinus carolus-bucarii C.P.Sm.
Lupinus chachas Ochoa ex C. P. Smith
Lupinus chamissonis Eschscholtz
Lupinus chiapensis Rose
Lupinus chihuahuensis S.Watson
Lupinus christianus C.P.Sm.
Lupinus chrysomelas Casar.
Lupinus clementinus Greene
Lupinus comatus Rydb.
Lupinus consentinii Walp.
Lupinus cymb-Aegressus C.P.Sm.
Lupinus dasyphyllus Greene
Lupinus davisianus C.P.Sm.
Lupinus debilis Eastw.
Lupinus decaschistus C.P.Sm.
Lupinus diaboli-septem C.P.Sm.
Lupinus dichrous Greene
Lupinus dispersus A.Heller
Lupinus dissimulans C.P.Sm.
Lupinus durangensis C.P.Sm.
Lupinus eatonanus C.P.Sm.
Lupinus equi-coeli C.P.Sm.
Lupinus equi-collis C.P.Sm.
Lupinus erectus L.F.Hend.
Lupinus erminens S.Watson
Lupinus ermineus S.Watson
Lupinus falcifer Nutt.
Lupinus falsoerectus C.P.Sm.
Lupinus falsoformosus C.P.Sm.
Lupinus falsograyi C.P.Sm.
Lupinus fieldii Rose ex J. F. Macbr.
Lupinus filicaulis C.P.Sm.
Lupinus finitus C.P.Sm.
Lupinus flavescens Rydb.
Lupinus foliosus Hook.
Lupinus foliosus Nutt.
Lupinus forskahlei Boiss.
Lupinus franciscanus Greene
Lupinus fraxinetorum Greene
Lupinus fruticosus Steud.
Lupinus fruticosus Dum.Cours.
Lupinus garcianus Bennett & Dunn
Lupinus geophilus Rose
Lupinus geraniophilus C.P.Sm.
Lupinus glabellus M.Martens & Galeotti
Lupinus graciliflorus C.P.Sm.
Lupinus gratus Greene
Lupinus gredensis Gand.
Lupinus guadalupensis Greene
Lupinus guadiloupensis Steud.
Lupinus guatimalensis auct.
Lupinus gussoneanus J.Agardh
Lupinus habrocomus Greene
Lupinus haudcytisoides C.P.Sm.
Lupinus helleri Greene
Lupinus hexaedrus E. Fourn.
Lupinus hintonii C.P.Sm.
Lupinus huigrensis Rose ex C.P.Sm.
Lupinus humicolus A.Nelson
Lupinus humifusus Benth.
Lupinus humilis Rose ex Pittier
Lupinus hyacinthinus Greene
Lupinus idoneus C.P.Sm.
Lupinus inamoenus Greene ex C.F.Baker
Lupinus indutus Greene ex C.F.Baker
Lupinus insignis Glaz. ex C. P. Smith
Lupinus integrifolius L.
Lupinus intergrifolius Desr.
Lupinus ione-grisetae C.P.Sm.
Lupinus ione-walkerae C.P.Sm.
Lupinus jamesonianus C.P.Sm.
Lupinus javanicus Burm.f.
Lupinus jorgensenanus C.P.Sm.
Lupinus jucundus Greene
Lupinus kellerrnanianus C.P.Sm.
Lupinus kyleanus C.P.Sm.
Lupinus labiatus Nutt.
Lupinus lacticolor Tamayo
Lupinus lacus-huntingtonii C.P.Sm.
Lupinus lacuum-trinitatum C.P.Sm.
Lupinus larsonanus C.P.Sm.
Lupinus lassenensis Eastw.
Lupinus latissimus Greene
Lupinus laxifolius A.Gray
Lupinus leptostachyus Greene
Lupinus lesueurii Standl.
Lupinus linearifolius Larrañaga
Lupinus lingulae C.P.Sm.
Lupinus longilabrum C.P.Sm.
Lupinus lorentzianus C.P.Sm.
Lupinus louise-bucariae C.P.Sm.
Lupinus louise-grisetae C.P.Sm.
Lupinus lucidus Benth. ex Loudon
Lupinus lyman-bensonii C.P.Sm.
Lupinus lysichitophilus C.P.Sm.
Lupinus macrocarpus Hook. & Arn.
Lupinus macrocarpus Torr.
Lupinus macrophyllus Benth.
Lupinus macrorhizos Georgi
Lupinus magnistipulatus Planchuelo & Dunn
Lupinus maissurianii Atabek. & Polukhina
Lupinus marcusianus C.P.Sm.
Lupinus mariae-josephae H.Pascual
Lupinus markleanus C.P.Sm.
Lupinus marschallianus Sweet
Lupinus mearnsii C.P.Sm.
Lupinus meli-campestris C.P.Sm.
Lupinus meridanus Moritz ex C. P. Smith
Lupinus mexiae C.P.Sm.
Lupinus micensis M.E.Jones
Lupinus micheneri Greene
Lupinus milleri J.Agardh
Lupinus minearanus C.P.Sm.
Lupinus minutissimus Tamayo
Lupinus molle A.Heller
Lupinus mollissifolius Davidson
Lupinus monettianus C.P.Sm.
Lupinus muellerianus C.P.Sm.
Lupinus multicincinnis C.P.Sm.
Lupinus neglectus Rose
Lupinus nemoralis Greene
Lupinus niger Wehmer
Lupinus noldekae Eastw.
Lupinus nutcanus Spreng.
Lupinus nutkatensis J.G.Cooper
Lupinus obtunsus C.P.Sm.
Lupinus octablomus C.P.Sm.
Lupinus opsianthus Amabekova & Maisuran
Lupinus pavonum C.P.Sm.
Lupinus pendeltonii A.Heller
Lupinus pendletonii A.Heller
Lupinus perconfertus C.P.Sm.
Lupinus perplexus C.P.Sm.
Lupinus philistaeus Boiss.
Lupinus pinus-contortae C.P.Sm.
Lupinus piperi B.L.Rob. ex Piper
Lupinus piperitus Davidson
Lupinus platanophilus M.E.Jones
Lupinus plebeius Greene ex C.F.Baker
Lupinus prato-lacuum C.P.Sm.
Lupinus prolifer Desr.
Lupinus propinquus Greene
Lupinus proteanus Eastw.
Lupinus psoraleoides Pollard
Lupinus pumviridis C.P.Sm.
Lupinus puroviridis C.P.Sm.
Lupinus purpurascens A.Heller
Lupinus pygmaeus Tamayo
Lupinus quercus-jugi C.P.Sm.
Lupinus quercuum C.P.Sm.
Lupinus rainierensis Eastw.
Lupinus regius Rudolph ex Torr. & A.Gray
Lupinus rhodanthus C.P.Sm.
Lupinus rickeri C.P.Sm.
Lupinus rivetianus C.P.Sm.
Lupinus rydbergii Blank.
Lupinus sabuli C.P.Sm.
Lupinus salicisocius C.P.Sm.
Lupinus salinensis C.P.Sm.
Lupinus sativus Gaterau
Lupinus scaposus Rydb.
Lupinus scheuberae Rydb.
Lupinus schickendantzii C.P.Sm.
Lupinus schiedeanus Steud.
Lupinus schumannii C.P.Sm.
Lupinus seclusus C.P.Sm.
Lupinus semiaequus C.P.Sm.
Lupinus semiverticillatus Desr.
Lupinus sergenti Tamayo ex Pittier
Lupinus sergentii Tamayo
Lupinus serradentum C.P.Sm.
Lupinus shrevei C.P.Sm.
Lupinus sierrae-zentae C.P.Sm.
Lupinus sileri S.Watson
Lupinus sinus-meyersii C.P. Sm.
Lupinus sparhawkianus C.P.Sm.
Lupinus spatulata Larrañaga
Lupinus speciosus Voss
Lupinus spruceanus C.P.Sm.
Lupinus standleyensis C.P.Sm.
Lupinus stationis C.P.Sm.
Lupinus stiveri Kellogg
Lupinus stoloniferus L.
Lupinus strigulosus Gand.
Lupinus subhirsutus Davidson
Lupinus subvolutus C.P.Sm.
Lupinus suksdorfii B.L. Rob. ex Piper
Lupinus summersianus C.P.Sm.
Lupinus sylvaticus Hemsl.
Lupinus thermis Gasp.
Lupinus thermus St.-Lag.
Lupinus tilcaricus C.P.Sm.
Lupinus timotensis Tamayo
Lupinus tricolor Greene
Lupinus tricolor G.Nicholson
Lupinus trifidus Torr. ex S.Watson
Lupinus tristis Sweet
Lupinus trochophyllus Hoffmanns.
Lupinus tuckeranus C.P. Sm.
Lupinus vaginans Benth.
Lupinus valdepallidus C.P.Sm.
Lupinus vandykeae Eastw.
Lupinus variegatus A.Heller
Lupinus variegatus Poir.
Lupinus varneranus C.P.Sm.
Lupinus vavilovii Atabekova & Maissurjan
Lupinus venustus Bailly
Lupinus violaceus A.Heller
Lupinus viridicalyx C.P.Sm.
Lupinus volcanicus Greene
Lupinus watsonii A.Heller
Lupinus westiana Small
Lupinus wolfianus C.P.Sm.
Lupinus yanlyensis C.P.Sm.
Lupinus yaruahensis C.P.Sm.
Hybrids
The following hybrids have been described:
Lupinus ×alpestris (A. Nelson) D.B. Dunn & J.M. Gillett
Lupinus ×hispanicoluteus W.Święcicki & W.K.Święcicki
Lupinus ×hybridus Lem.
Lupinus ×insignis Lem.
Lupinus ×regalis (auct.) Bergmans—rainbow lupin (Lupinus arboreus × Lupinus polyphyllus)
Lupinus ×versicolor Caball.
Etymology
While some sources believe the origin of the name to be in doubt , the Collins Dictionary definition asserts that the word is 14th century in origin, from the Latin lupīnus "wolfish" from lupus "wolf" as it was believed that the plant ravenously exhausted the soil.
But a more likely explanation is that lupinus meant that the plants were as dangerous to livestock as wolves, because the alkaloid poisons of Lupines can sicken or kill grazing animals, especially sheep. Farmers have known since ancient Rome that lupines improve soil by adding nitrogen and loosening compacted earth with their strong root systems, so the Collins explanation is improbable.
Ecology
Certain species, such as the yellow bush lupin (L. arboreus), are considered invasive weeds when they appear outside their native ranges. In New Zealand, lupines are viewed as invasive and a severe threat in some cases. L. polyphyllus has escaped into the wild and grows in large numbers along main roads and streams on the South Island. A similar spread of the species has occurred in Sweden, Finland and Norway after the non-native species was first deliberately planted in the landscaping along the main roads. Lupins have been planted in some parts of Australia with a considerably cooler climate, particularly in rural Victoria and New South Wales.
Lupins are important larval food plants for many lepidopterans (butterflies and moths). These include:
Iraricia icarioides missionensis (Mission blue butterfly), larvae limited to Lupinus
Callophrys irus (frosted elfin), recorded on L. perennis
Erynnis persius (Persius duskywing)
†Glaucopsyche xerces (Xerces blue)
Glaucopsyche lygdamus (silvery blue)
Plebejus melissa samuelis (Karner blue)
Erynnis persius persius (eastern Persius duskywing)
Schinia sueta, larvae limited to Lupinus
Cultivation
Lupinus polyphyllus, the garden lupin, and Lupinus arboreus, the tree lupin, are popular ornamental plants in gardens, and are the source of numerous hybrids and cultivars in a wide range of colours, including bicolors. As legumes, lupins are good companion plants in gardens, increasing the soil nitrogen for vegetables and other plants. As well as growing in the ground, lupins can do well in pots on balconies or patios.
Agriculture
Like other legumes, lupines can fix nitrogen from the atmosphere into ammonia via a rhizobium–root nodule symbiosis, fertilizing the soil for other plants. This adaptation allows lupins to be tolerant of infertile soils and capable of pioneering change in barren and poor-quality soils. The genus Lupinus is nodulated by Bradyrhizobium soil bacteria.
In the early 20th century, German scientists attempted to cultivate a sweet variety of lupin lacking the bitter taste, making it more suitable for both human and animal consumption.
Many annual species of lupins are used in agriculture and most of them have Mediterranean origin.
While originally cultivated as a green manure or forage, lupins are increasingly grown for their seeds, which can be used as an alternative to soybeans. Sweet (low alkaloid) lupins are highly regarded as a stock feed, particularly for ruminants, but also for pigs and poultry and more recently as an ingredient in aqua-feeds. Three Mediterranean species of lupin, blue (narrow-leafed) lupin, white lupin, and yellow lupin, are widely cultivated for livestock and poultry feed.
The market for lupin seeds for human food is currently small, but researchers believe it has great potential. Lupin seeds are considered "superior" to soybeans in certain applications and evidence is increasing for their potential health benefits. They contain similar protein to soybean, but less fat. As a food source, they are gluten-free and high in dietary fibre, amino acids, and antioxidants, and they are considered to be prebiotic.
About 85% of the world's lupin seeds are grown in Western Australia.
Toxicity
Some lupins contain certain secondary compounds, including isoflavones and toxic alkaloids, such as lupinine, anagyrine and sparteine. With early detection, these can be removed through processing, although lupins containing these elements are not usually selected for food-grade products.
A risk of lupin allergy exists in patients allergic to peanuts. Most lupin reactions reported have been in people with peanut allergy. Because of the cross-allergenicity of peanut and lupin, the European Commission, as of 2006, has required that food labels indicate the presence of "lupin and products thereof" in food.
Lupin plants can be colonized by the fungus Diaporthe toxica which can cause a mycotoxicosis known as lupinosis when ingested by grazing animals.
Uses
The legume seeds of lupins, commonly called lupin beans, were popular with the Romans, who cultivated the plants throughout the Roman Empire where the lupin is still known in extant Romance languages by names such as .
Seeds of various species of lupins have been used as a food for over 3,000 years around the Mediterranean and for as long as 6,000 years in the Andes. Lupins were also used by many Native American peoples of North America such as the Yavapai. The Andean lupin or (Lupinus mutabilis) was a widespread food in the Incan Empire; but they have never been accorded the same status as soybeans, dry peas and other pulse crops. The pearl lupin of the Andean highlands of South America, L. mutabilis, known locally as or , was extensively cultivated, but no conscious genetic improvement other than to select for larger and water-permeable seeds seems to have been made. Users soaked the seed in running water to remove most of the bitter alkaloids and then cooked or toasted the seeds to make them edible, or else boiled and dried them to make , reported as a pre-Columbian practice in . Spanish domination led to a change in the eating habits of the indigenous peoples, and only recently (late 20th century onward) has interest in using lupins as a food been renewed.
Lupins can be used to make a variety of foods both sweet and savoury, including everyday meals, traditional fermented foods, baked foods, and sauces. The European white lupin (L. albus) beans are commonly sold in a salty solution in jars (like olives and pickles) and can be eaten with or without the skin. Lupini dishes are most commonly found in Europe, especially in Portugal, Spain, Greece, and Italy. They are also common in Brazil and Egypt. In Egypt, the lupin is known in Arabic as , and is a popular street snack after being treated with several soakings of water, and then brined. In Portugal, Spain, and the Spanish Harlem district of New York, they are consumed with beer and wine. In Lebanon, Palestine, Israel, Jordan, and Syria the salty and chilled lupini beans are called turmus (in , ) and are served as part of an apéritif or a snack. Other species, such as L. albus (white lupin), L. angustifolius (narrow-leafed lupin), and L. hirsutus (blue lupin) also have edible seeds.
Culture
Consumed throughout the Mediterranean region and the Andean mountains, lupins were eaten by the early Egyptian and pre-Incan people and were known to Roman agriculturalists for their ability to improve the fertility of soils.
In the late 18th century, lupins were introduced into northern Europe as a means of improving soil quality, and by the 1860s, the garden yellow lupin was seen across the sandy soils of the Baltic coastal plain.
The successful development of lupin varieties with the necessary "sweet gene" paved the way for the greater adoption of lupins across Europe and later Australia.
Further work carried out by the Western Australian Department of Agriculture and Food during the 1950s and '60s has led to more sweet lupin crops being produced in Western Australia now than anywhere else in the world.
Bluebonnets, including the Texas bluebonnet (L. texensis), are the state flowers of Texas.
| Biology and health sciences | Fabales | Plants |
305742 | https://en.wikipedia.org/wiki/Squamata | Squamata | Squamata (, Latin squamatus, 'scaly, having scales') is the largest order of reptiles, comprising lizards and snakes. With over 12,162 species, it is also the second-largest order of extant (living) vertebrates, after the perciform fish. Squamates are distinguished by their skins, which bear horny scales or shields, and must periodically engage in molting. They also possess movable quadrate bones, making possible movement of the upper jaw relative to the neurocranium. This is particularly visible in snakes, which are able to open their mouths very widely to accommodate comparatively large prey. Squamates are the most variably sized living reptiles, ranging from the dwarf gecko (Sphaerodactylus ariasae) to the reticulated python (Malayopython reticulatus). The now-extinct mosasaurs reached lengths over .
Among other reptiles, squamates are most closely related to the tuatara, the last surviving member of the once diverse Rhynchocephalia, with both groups being placed in the clade Lepidosauria.
Evolution
Squamates are a monophyletic sister group to the rhynchocephalians, members of the order Rhynchocephalia. The only surviving member of the Rhynchocephalia is the tuatara. Squamata and Rhynchocephalia form the subclass Lepidosauria, which is the sister group to the Archosauria, the clade that contains crocodiles and birds, and their extinct relatives. Fossils of rhynchocephalians first appear in the Early Triassic, meaning that the lineage leading to squamates must have also existed at the time.
A study in 2018 found that Megachirella, an extinct genus of lepidosaurs that lived about 240 million years ago during the Middle Triassic, was a stem-squamate, making it the oldest known squamate. The phylogenetic analysis was conducted by performing high-resolution microfocus X-ray computed tomography (micro-CT) scans on the fossil specimen of Megachirella to gather detailed data about its anatomy. These data were then compared with a phylogenetic dataset combining the morphological and molecular data of 129 extant and extinct reptilian taxa. The comparison revealed Megachirella had certain features that are unique to squamates. The study also found that geckos are the earliest crown group squamates, not iguanians. However, a 2021 study found the genus to be a lepidosaur of uncertain position, in a polytomy with Squamata and Rhynchocephalia.
In 2022, the extinct genus Cryptovaranoides was described from the Late Triassic (Rhaetian age) of England as a highly derived squamate belonging to the group Anguimorpha, which contains many extant lineages such as monitor lizards, beaded lizards and anguids. The presence of an essentially modern crown group squamate so far back in time was unexpected, as their diversification was previously thought to have occurred during the Jurassic and Cretaceous. A 2023 study found that Cryptovaranoides most likely represents an archosauromorph with no apparent squamate affinities, though the original describers maintained their original conclusion that this taxon represents a squamate. The oldest unambiguous fossils of Squamata date to the Bathonian age of the Middle Jurassic of the Northern Hemisphere, with the first appearance of many modern groups, including snakes, during this period.
Scientists believe crown group squamates probably originated in the Early Jurassic based on the fossil record, with the oldest unambiguous fossils of squamates dating to the Middle Jurassic. Squamate morphological and ecological diversity substantially increased over the course of the Cretaceous, including the appeance of groups like iguanians and varanoids, and true snakes. Polyglyphanodontia, an extinct clade of lizards, and mosasaurs, a group of predatory marine lizards that grew to enormous sizes, also appeared in the Cretaceous. Squamates suffered a mass extinction at the Cretaceous–Paleogene (K–Pg) boundary, which wiped out polyglyphanodontians, mosasaurs, and many other distinct lineages.
The relationships of squamates are debatable. Although many of the groups originally recognized on the basis of morphology are still accepted, understanding of their relationships to each other has changed radically as a result of studying their genomes. Iguanians were long thought to be the earliest crown group squamates based on morphological data, but genetic data suggest that geckos are the earliest crown group squamates. Iguanians are now united with snakes and anguimorphs in a clade called Toxicofera. Genetic data also suggest that the various limbless groups – snakes, amphisbaenians, and dibamids – are unrelated, and instead arose independently from lizards.
Reproduction
The male members of the group Squamata have hemipenes, which are usually held inverted within their bodies, and are everted for reproduction via erectile tissue like that in the mammalian penis. Only one is used at a time, and some evidence indicates that males alternate use between copulations. The hemipenis has a variety of shapes, depending on the species. Often it bears spines or hooks, to anchor the male within the female. Some species even have forked hemipenes (each hemipenis has two tips). Due to being everted and inverted, hemipenes do not have a completely enclosed channel for the conduction of sperm, but rather a seminal groove that seals as the erectile tissue expands. This is also the only reptile group in which both viviparous and ovoviviparous species are found, as well as the usual oviparous reptiles. The eggs in oviparous species have a parchment-like shell. The only exception is found in blind lizards and three families of geckos (Gekkonidae, Phyllodactylidae and Sphaerodactylidae), where many lay rigid and calcified eggs. Some species, such as the Komodo dragon, can reproduce asexually through parthenogenesis.
Studies have been conducted on how sexual selection manifests itself in snakes and lizards. Snakes use a variety of tactics in acquiring mates. Ritual combat between males for the females with which they want to mate includes topping, a behavior exhibited by most viperids, in which one male twists around the vertically elevated fore body of his opponent and forcing it downward. Neck biting commonly occurs while the snakes are entwined.
Facultative parthenogenesis
Parthenogenesis is a natural form of reproduction in which the growth and development of embryos occur without fertilization. Agkistrodon contortrix (copperhead snake) and Agkistrodon piscivorus (cottonmouth snake) can reproduce by facultative parthenogenesis; they are capable of switching from a sexual mode of reproduction to an asexual mode. The type of parthenogenesis that likely occurs is automixis with terminal fusion (see figure), a process in which two terminal products from the same meiosis fuse to form a diploid zygote. This process leads to genome-wide homozygosity, expression of deleterious recessive alleles, and often to developmental abnormalities. Both captive-born and wild-born A. contortrix and A. piscivorus appear to be capable of this form of parthenogenesis.
Reproduction in squamate reptiles is ordinarily sexual, with males having a ZZ pair of sex-determining chromosomes, and females a ZW pair. However, the Colombian rainbow boa, Epicrates maurus, can also reproduce by facultative parthenogenesis, resulting in production of WW female progeny. The WW females are likely produced by terminal automixis.
Inbreeding avoidance
When female sand lizards mate with two or more males, sperm competition within the female's reproductive tract may occur. Active selection of sperm by females appears to occur in a manner that enhances female fitness. On the basis of this selective process, the sperm of males that are more distantly related to the female are preferentially used for fertilization, rather than the sperm of close relatives. This preference may enhance the fitness of progeny by reducing inbreeding depression.
Evolution of venom
Recent research suggests that the evolutionary origin of venom may exist deep in the squamate phylogeny, with 60% of squamates placed in this hypothetical group called Toxicofera. Venom has been known in the clades Caenophidia, Anguimorpha, and Iguania, and has been shown to have evolved a single time along these lineages before the three groups diverged, because all lineages share nine common toxins. The fossil record shows the divergence between anguimorphs, iguanians, and advanced snakes dates back roughly 200 million years ago (Mya) to the Late Triassic/Early Jurassic, but the only good fossil evidence is from the Middle Jurassic.
Snake venom has been shown to have evolved via a process by which a gene encoding for a normal body protein, typically one involved in key regulatory processes or bioactivity, is duplicated, and the copy is selectively expressed in the venom gland. Previous literature hypothesized that venoms were modifications of salivary or pancreatic proteins, but different toxins have been found to have been recruited from numerous different protein bodies and are as diverse as their functions.
Natural selection has driven the origination and diversification of the toxins to counter the defenses of their prey. Once toxins have been recruited into the venom proteome, they form large, multigene families and evolve via the birth-and-death model of protein evolution, which leads to a diversification of toxins that allows the ambush predators the ability to attack a wide range of prey. The rapid evolution and diversification is thought to be the result of a predator–prey evolutionary arms race, where both are adapting to counter the other.
Humans and squamates
Bites and fatalities
An estimated 125,000 people a year die from venomous snake bites. In the US alone, more than 8,000 venomous snake bites are reported each year, but only one in 50 million people (five or six fatalities per year in the USA) will die from venomous snake bites.
Lizard bites, unlike venomous snake bites, are usually not fatal. The Komodo dragon has been known to kill people due to its size, and recent studies show it may have a passive envenomation system. Recent studies also show that the close relatives of the Komodo, the monitor lizards, all have a similar envenomation system, but the toxicity of the bites is relatively low to humans. The Gila monster and beaded lizards of North and Central America are venomous, but not deadly to humans.
Conservation
Though they survived the Cretaceous–Paleogene extinction event, many squamate species are now endangered due to habitat loss, hunting and poaching, illegal wildlife trading, alien species being introduced to their habitats (which puts native creatures at risk through competition, disease, and predation), and other anthropogenic causes. Because of this, some squamate species have recently become extinct, with Africa having the most extinct species. Breeding programs and wildlife parks, though, are trying to save many endangered reptiles from extinction. Zoos, private hobbyists, and breeders help educate people about the importance of snakes and lizards.
Classification and phylogeny
Historically, the order Squamata has been divided into three suborders:
Lacertilia, the lizards
Serpentes, the snakes (see also Ophidia)
Amphisbaenia, the worm lizards
Of these, the lizards form a paraphyletic group, since the "lizards" are found in several distinct lineages, with snakes and amphisbaenians recovered as monophyletic groups nested within. Although studies of squamate relationships using molecular biology have found different relationships between some squamata lineagaes, all recent molecular studies suggest that the venomous groups are united in a venom clade. Named Toxicofera, it encompasses a majority (nearly 60%) of squamate species and includes Serpentes (snakes), Iguania (agamids, chameleons, iguanids, etc.), and Anguimorpha (monitor lizards, Gila monster, glass lizards, etc.).
One example of a modern classification of the squamates is shown below.
List of extant families
The over 10,900 extant squamates are divided into 67 families.
| Biology and health sciences | Lizards and other Squamata | Animals |
305746 | https://en.wikipedia.org/wiki/Sauria | Sauria | Sauria is the clade of diapsids containing the most recent common ancestor of Archosauria (which includes crocodilians and birds) and Lepidosauria (which includes squamates and the tuatara), and all its descendants. Since most molecular phylogenies recover turtles as more closely related to archosaurs than to lepidosaurs as part of Archelosauria, Sauria can be considered the crown group of diapsids, or reptiles in general. Depending on the systematics, Sauria includes all modern reptiles or most of them (including birds, a type of archosaur) as well as various extinct groups.
Sauria lies within the larger total group Sauropsida, which also contains various stem-reptiles which are more closely related to reptiles than to mammals. Prior to its modern usage, "Sauria" was used as a name for the suborder occupied by lizards, which before 1800 were considered crocodilians.
Systematics
Sauria was historically used as a partial equivalent for Squamata (which contains lizards and snakes). The redefinition to cover the last common ancestor of archosaurs and lepidosaurs was the result of papers by Jacques A. Gauthier and colleagues in the 1980s.
Genomic studies and comprehensive studies in the fossil record suggest that turtles are closely related to archosaurs as part of Sauria, and not to the non-saurian parareptiles as previously thought. In a 2018 cladistic analysis, Pantestudines (turtles and close relatives) were placed within Diapsida but outside of Sauria.
Synapomorphies
The synapomorphies or characters that unite the clade Sauria also help them be distinguished from stem-saurians in Diapsida or stem-reptiles in clade Sauropsida in the following categories based on the following regions of the body.
Cephalad Region
Dorsal origin of temporal musculature
Loss of caniniform region in maxillary tooth row
External nares close to the midline
Postparietal absent
Squamosal mainly restricted to top of skull
The occipital flange of the squamosal is little exposed on the occiput
Anterior process of squamosal narrow
Quadrate exposed laterally
Unossified dorsal process of stapes
Stapes slender
Trunk Region
Sacral ribs oriented laterally
Ontogenetic fusion of caudal ribs
Trunk ribs mostly single headed
Pectoral Region
Cleithrum absent
Pelvic Region
Modified ilium
Limb Region
Tubular bone lost
Entepicondylar foramen absent
Radius as long as ulna
Small proximal carpals and tarsal
Fifth distal tarsal absent
Short and stout fifth or hooked metatarsal
Perforating foramen of manus lost
However, some of these characters might be lost or modified in several lineages, particularly among birds and turtles; it is best to see these characters as the ancestral features that were present in the ancestral saurian.
Phylogeny
The cladogram shown below follows the most likely result found by an analysis of turtle relationships using both fossil and genetic evidence by M.S. Lee, in 2013. This study found Eunotosaurus, usually regarded as a turtle relative, to be only very distantly related to turtles in the clade Parareptilia.
The cladogram below follows the most likely result found by another analysis of turtle relationships, this one using only fossil evidence, published by Rainer Schoch and Hans-Dieter Sues in 2015. This study found Eunotosaurus to be an actual early stem-turtle, though other versions of the analysis found weak support for it as a parareptile.
The cladogram below follows the analysis of Li et al. (2018). It places turtles within Diapsida but outside of Sauria (the Lepidosauromorpha + Archosauromorpha clade).
The following cladogram was found by Simões et al. (2022):
| Biology and health sciences | Reptiles: General | Animals |
305854 | https://en.wikipedia.org/wiki/Text%20messaging | Text messaging | Text messaging, or simply texting, is the act of composing and sending electronic messages, typically consisting of alphabetic and numeric characters, between two or more users of mobile phones, tablet computers, smartwatches, desktops/laptops, or another type of compatible computer. Text messages may be sent over a cellular network or may also be sent via satellite or Internet connection.
The term originally referred to messages sent using the Short Message Service (SMS) on mobile devices. It has grown beyond alphanumeric text to include multimedia messages using the Multimedia Messaging Service (MMS) and Rich Communication Services (RCS), which can contain digital images, videos, and sound content, as well as ideograms known as emoji (happy faces, sad faces, and other icons), and on various instant messaging apps. Text messaging has been an extremely popular medium of communication since the turn of the century and has also influenced changes in society.
Overview
Text messages are used for personal, family, business, and social purposes. Governmental and non-governmental organizations use text messaging for communication between colleagues. In the 2010s, the sending of short informal messages became an accepted part of many cultures, as happened earlier with emailing. This makes texting a quick and easy way to communicate with friends, family, and colleagues, including in contexts where a call would be impolite or inappropriate (e.g., calling very late at night or when one knows the other person is busy with family or work activities). Like e-mail and voicemail, and unlike calls (in which the caller hopes to speak directly with the recipient), texting does not require the caller and recipient to both be free at the same moment; this permits communication even between busy individuals. Text messages can also be used to interact with automated systems, for example, to order products or services from e-commerce websites or to participate in online contests. Advertisers and service providers use direct text marketing to send messages to mobile users about promotions, payment due dates, and other notifications instead of using postal mail, email, or voicemail.
Terminology
The service is referred to by different colloquialisms depending on the region. It may simply be referred to as a "text" in North America, the United Kingdom, Australia, New Zealand, and the Philippines; an "SMS" in most of mainland Europe; or an "MMS" or "SMS" in the Middle East, Africa, and Asia. The sender of a text message is commonly referred to as a "texter".
History
The electrical telegraph systems, developed in the early 19th century, used electrical signals to send text messages. In the late 19th century, wireless telegraphy was developed using radio waves.
In 1933, the German Reichspost (Reich postal service) introduced the first "telex" service.
The University of Hawaii began using radio to send digital information as early as 1971, using ALOHAnet. Friedhelm Hillebrand conceptualised SMS in 1984 while working for Deutsche Telekom. Sitting at a typewriter at home, Hillebrand typed out random sentences and counted every letter, number, punctuation mark, and space. Almost every time, the messages contained fewer than 160 characters, thus giving the basis for the limit one could type via text messaging. With Bernard Ghillebaert of France Télécom, he developed a proposal for the GSM (Groupe Spécial Mobile) meeting in February 1985 in Oslo. The first technical solution evolved in a GSM subgroup under the leadership of Finn Trosby. It was further developed under the leadership of Kevin Holley and Ian Harris (see Short Message Service). SMS forms an integral part of Signalling System No. 7 (SS7). Under SS7, it is a "state" with 160 characters of data, coded in the ITU-T "T.56" text format, that has a "sequence lead in" to determine different language codes and may have special character codes that permit, for example, sending simple graphs as text. This was part of ISDN (Integrated Services Digital Network), and since GSM is based on this, it made its way to the mobile phone. Messages could be sent and received on ISDN phones, and these can send SMS to any GSM phone. The possibility of doing something is one thing; implementing it is another, but systems existed in 1988 that sent SMS messages to mobile phones (compare ND-NOTIS).
SMS messaging was used for the first time on 3 December 1992, when Neil Papworth, a 22-year-old test engineer, used a computer to send the text message "Merry Christmas" via the Vodafone network to the phone of Richard Jarvis, who was at a party in Newbury, Berkshire celebrating the event. Papworth later said "it didn't feel momentous at all". Modern SMS text messaging is usually sent from one mobile phone to another. Finnish Radiolinja became the first network to offer a commercial person-to-person SMS text messaging service in 1994. When Radiolinja's domestic competitor, Telecom Finland (now part of TeliaSonera), also launched SMS text messaging in 1995 and the two networks offered cross-network SMS functionality, Finland became the first nation where SMS text messaging was offered on a competitive as well as a commercial basis. GSM was allowed in the United States, but the radio frequencies were blocked and awarded to US "Carriers" to use US technology, which limited development of mobile messaging services in the US. The GSM in the US had to use a frequency allocated for private communication services (PCS) – what the ITU frequency régime had blocked for DECT (Digital Enhanced Cordless Telecommunications) – a 1,000-foot range picocell, but it survived. American Personal Communications (APC), the first GSM carrier in America, provided the first text-messaging service in the United States. Sprint Telecommunications Venture, a partnership of Sprint Corp. and three large cable-TV companies, owned 49 percent of APC. The Sprint venture was the largest single buyer at a government-run spectrum auction that raised $7.7 billion in 2005 for PCS licenses. APC operated under the brand name Sprint Spectrum and launched its service on 15 November 1995, in Washington, D.C., and Baltimore, Maryland. Vice President Al Gore in Washington, D.C., made the initial phone call to launch the network, calling Mayor Kurt Schmoke in Baltimore.
Initial growth of text messaging worldwide was slow, with customers in 1995 sending on average only 0,4 messages per GSM customer per month. One factor in the slow take-up of SMS was that operators were slow to set up charging systems, especially for prepaid subscribers, and to eliminate billing fraud, which was possible by changing SMSC settings on individual handsets to use the SMSCs of other operators. Over time, this issue was eliminated by switch billing instead of billing at the SMSC and by new features within SMSCs that allowed the blocking of foreign mobile users sending messages through them. SMS is available on a wide range of networks, including 3G networks. However, not all text-messaging systems use SMS; some notable alternate implementations of the concept include J-Phone's SkyMail and NTT Docomo's Short Mail, both in Japan. E-mail messaging from phones, as popularized by NTT Docomo's i-mode and the RIM BlackBerry, also typically use standard mail protocols such as SMTP over TCP/IP. , text messaging was the most widely used mobile data service, with 74% of all mobile phone users worldwide, or 2.4 billion out of 3.3 billion phone subscribers, being active users of the Short Message Service at the end of 2007. In countries such as Finland, Sweden, and Norway, over 85% of the population used SMS. The European average was about 80%, and North America was rapidly catching up, with over 60% active users of SMS . The largest average usage of the service by mobile phone subscribers occurs in the Philippines, with an average of 27 texts sent per day per subscriber.
Uses
Text messaging is most often used between private mobile phone users, as a substitute for voice calls in situations where voice communication is impossible or undesirable (e.g., during a school class or a work meeting). Texting is also used to communicate very brief messages, such as informing someone that you will be late or reminding a friend or colleague about a meeting. As with e-mail, informality and brevity have become an accepted part of text messaging. Some text messages such as SMS can also be used for the remote control of home appliances. It is widely used in domotics systems. Some amateurs have also built their own systems to control (some of) their appliances via SMS. A Flash SMS is a type of text message that appears directly on the main screen without user interaction and is not automatically stored in the inbox. It can be useful in cases such as an emergency (e.g., fire alarm) or confidentiality (e.g., one-time password).
SMS has historically been particularly popular in Europe, Asia (excluding Japan; see below), the United States, Australia, and New Zealand, while also gaining influence in Africa. Popularity has grown to a sufficient extent that the term texting (used as a verb meaning the act of mobile phone users sending short messages back and forth) has entered the common lexicon. In 2012, young Asians considered SMS as the most popular mobile phone application. In the same year, 50 percent of American teens send 50 text messages or more per day, making it their most frequent form of communication. In 2004 in China, SMS was very popular and brought service providers significant profit (18 billion short messages were sent in 2001).
It has been a very influential and powerful tool in the Philippines, where in 2008 the average user sent 10–12 text messages a day. The same year, the Philippines alone sent on average over 1 billion text messages a day, more than the annual average SMS volume of the countries in Europe, and even China and India. SMS saw hugely popular in India, where youngsters often exchanged many text messages, and companies provide alerts, infotainment, news, cricket scores updates, railway/airline booking, mobile billing, and banking services on SMS.
Similarly, in 2008, text messaging played a primary role in the implication of former Detroit Mayor Kwame Kilpatrick in an SMS sex scandal. Short messages are particularly popular among young urbanites. In many markets, the service is comparatively cheap. For example, in Australia, a message typically costs between A$0.20 and $0.25 to send (some prepaid services charge $0.01 between their own phones), compared with a voice call, which costs somewhere between $0.40 and $2.00 per minute (commonly charged in half-minute blocks). The service is enormously profitable to the service providers. At a typical length of only 190 bytes (including protocol overhead), more than 350 of these messages per minute can be transmitted at the same data rate as a usual voice call (9 kbit/s). There are also free SMS services available, which are often sponsored, that allow sending and receiving SMS from a PC connected to the Internet. Mobile service providers in New Zealand, such as One NZ and Spark New Zealand, provided up to 2000 SMS messages for NZ$10 per month. Users on these plans sent on average 1500 SMS messages every month. Text messaging became so popular that advertising agencies and advertisers jumped into the text messaging business. Services that provide bulk text message sending are also becoming a popular way for clubs, associations, and advertisers to reach a group of opt-in subscribers quickly.
In 2013, research suggested that Internet-based mobile messaging would grow to equal the popularity of SMS by the end of 2013, with nearly 10 trillion messages being sent through each technology. Services such as Facebook Messenger/WhatsApp, Signal (software), Snapchat, Telegram (software), Viber have led to a decline in the use of SMS in parts of the world. A survey conducted by MetrixLabs showed that 63% of Baby Boomers, 63% of Generation X, and 67% of Generation Y said that they used instant messengers in place of texting. A Facebook survey showed that 65% of people surveyed thought that messaging applications made group messaging easier.
Applications
Microblogging
Of many texting trends, a system known as microblogging has surfaced, which consists of a miniaturized blog, inspired mainly by people's tendency to jot down informal thoughts and post them online. They consist of websites like X (formerly Twitter) and its Chinese equivalent Weibo (微博). As of 2016, 21% of all American adults used Twitter. As of 2017, Weibo had 340 million active users.
Emergency services
In some countries, text messages can be used to contact emergency services. In the UK, text messages can be used to call emergency services only after registering with the emergency SMS service. This service is primarily aimed at people who, because of disability, are unable to make a voice call. It has recently been promoted as a means for walkers and climbers to call emergency services from areas where a voice call is not possible due to low signal strength.
In the US, there is a move to require both traditional operators and over-the-top messaging providers to support texting to 911. In Asia, SMS is used for tsunami warnings and in Europe, SMS is used to inform individuals of imminent disasters. Since the location of a handset is known, systems can alert everyone in an area that the events have made impossible to pass through e.g. an avalanche. A similar system, known as Emergency Alert, is used in Australia to notify the public of impending disasters through both SMS and landline phone calls. These messages can be sent based on either the location of the phone or the address to which the handset is registered.
In the early 2020s, device manufacturers have begun to integrate satellite messaging connectivity and satellite emergency services into conventional mobile phones for use in remote regions, where there is no reliable terrestrial cellular network.
Reminders of medical appointments
SMS messages are used in some countries as reminders of medical appointments. Missed outpatient clinic appointments cost the National Health Service (England) more than £600 million ($980 million) a year. SMS messages are thought to be more cost-effective, swifter to deliver, and more likely to receive a faster response than letters. A 2012 study by Sims and colleagues examined the outcomes of 24,709 outpatient appointments scheduled in mental health services in South-East London. The study found that SMS message reminders could reduce the number of missed psychiatric appointments by 25–28%, representing a potential national yearly saving of over £150 million.
Because of the COVID-19 pandemic, medical facilities in the United States are using text messaging to coordinate the appointment process, including reminders, cancellations, and safe check-in. US-based cloud radiology information system vendor AbbaDox includes this in their patient engagement services.
Commercial uses
Short codes
Short codes are special telephone numbers, shorter than full telephone numbers, that can be used to address SMS and MMS messages from mobile phones or fixed phones. There are two types of short codes: dialling and messaging.
Text messaging gateway providers
SMS gateway providers facilitate the SMS traffic between businesses and mobile subscribers, being mainly responsible for carrying mission-critical messages, SMS for enterprises, content delivery and entertainment services involving SMS, e.g., TV voting. Considering SMS messaging performance and cost, as well as the level of text messaging services, SMS gateway providers can be classified as resellers of the text messaging capability of another provider's SMSC or offering the text messaging capability as an operator of their own SMSC with SS7. SMS messaging gateway providers can provide gateway-to-mobile (Mobile Terminated–MT) services. Some suppliers can also supply mobile-to-gateway (text-in or Mobile Originated/MO services). Many operate text-in services on short codes or mobile number ranges, whereas others use lower-cost geographic text-in numbers.
Premium content
SMS has been widely used for delivering digital content, such as news alerts, financial information, pictures, GIFs, logos and ringtones. Such messages are also known as premium-rated short messages (PSMS). The subscribers are charged extra for receiving this premium content, and the amount is typically divided between the mobile network operator and the value added service provider (VASP), either through revenue share or a fixed transport fee. Services like 82ASK and Any Question Answered have used the PSMS model to enable rapid response to mobile consumers' questions, using on-call teams of experts and researchers. In November 2013, amidst complaints about unsolicited charges on bills, major mobile carriers in the US agreed to stop billing for PSMS in 45 states, effectively ending its use in the United States.
Outside the United States, premium short messages have been used for "real-world" services. For example, some vending machines now allow payment by sending a premium-rated short message, so that the cost of the item bought is added to the user's phone bill or subtracted from the user's prepaid credits. Recently, premium messaging companies have come under fire from consumer groups due to a large number of consumers racking up huge phone bills. A new type of free-premium or hybrid-premium content has emerged with the launch of text-service websites. These sites allow registered users to receive free text messages when items they are interested in go on sale, or when new items are introduced. An alternative to inbound SMS is based on long numbers (international mobile number format, e.g., +44 7624 805000, or geographic numbers that can handle voice and SMS, e.g., 01133203040), which can be used in place of short codes or premium-rated short messages for SMS reception in several applications, such as TV voting, product promotions and campaigns. Long numbers are internationally available, as well as enabling businesses to have their own number, rather than short codes, which are usually shared across a lot of brands. Additionally, long numbers are non-premium inbound numbers.
In workplaces
The use of text messaging for workplace purposes grew significantly during the mid-2000s. As companies seek competitive advantages, many employees used new technology, collaborative applications, and real-time messaging such as SMS, instant messaging, and mobile communications to connect with teammates and customers. Some practical uses of text messaging include the use of SMS for confirming delivery or other tasks, for instant communication between a service provider and a client (e.g., a payment card company and a consumer), and for sending alerts. Several universities have implemented a system of texting students and faculties campus alerts. One such example is Penn State.
As text messaging has proliferated in business, so too have regulations governing its use. One regulation specifically governing the use of text messaging in financial-services firms engaged in stocks, equities, and securities trading is Regulatory Notice 07-59, Supervision of Electronic Communications, December 2007, issued to member firms by the Financial Industry Regulatory Authority (FINRA). In Regulatory Notice 07-59, FINRA noted that "electronic communications", "e-mail", and "electronic correspondence" may be used interchangeably and can include such forms of electronic messaging as instant messaging and text messaging. Industry has had to develop new technology to allow companies to archive their employees' text messages.
Security, confidentiality, reliability, and speed of SMS are among the most important guarantees industries such as financial services, energy and commodities trading, health care and enterprises demand in their mission-critical procedures. One way to guarantee such a quality of text messaging lies in introducing SLAs (Service Level Agreement), which are common in IT contracts. By providing measurable SLAs, corporations can define reliability parameters and set up a high quality of their services. Just one of many SMS applications that have proven highly popular and successful in the financial services industry is mobile receipts. In January 2009, Mobile Marketing Association (MMA) published the Mobile Banking Overview for financial institutions in which it discussed the advantages and disadvantages of mobile channel platforms such as Short Message Services (SMS), Mobile Web, Mobile Client Applications, SMS with Mobile Web and Secure SMS.
Mobile interaction services are an alternative way of using SMS in business communications with greater certainty. Typical business-to-business applications are telematics and Machine-to-Machine, in which two applications automatically communicate with each other. Incident alerts are also common, and staff communications are also another use for B2B scenarios. Businesses can use SMS for time-critical alerts, updates, and reminders, mobile campaigns, content and entertainment applications. Mobile interaction can also be used for consumer-to-business interactions, such as media voting and competitions, and consumer-to-consumer interaction, for example, with mobile social networking, chatting and dating.
Text messaging is widely used in business settings; as well, it is used in many civil service and non-governmental organization workplaces. The U.S. And Canadian civil service both adopted BlackBerry smartphones in the 2000s.
Group texts
Group texts involve more than two users. They are often used when it is helpful to message many people at once, such as inviting multiple people to an event or arranging groups. They are also used in business for marketing and other customer notifications as well as intracompany communication.
Group texts are often sent as MMS messages and therefore require an internet connection to send instead of using the sender's text messaging plan.
Online SMS services
There are a growing number of websites that allow users to send free SMS messages online. Some websites provide free SMS for promoting premium business packages.
Worldwide use
Europe
In 2003, Europe followed next behind Asia in terms of the popularity of the use of SMS. That year, an average of 16 billion messages were sent each month. Users in Spain sent a little more than fifty messages per month on average in 2003. In Italy, Germany and the United Kingdom, the figure was around 35–40 SMS messages per month. In each of these countries, the cost of sending an SMS message varied from €0.04–0.23, depending on the payment plan (with many contractual plans including all or several texts for free). In the United Kingdom, text messages are charged between £0.05–0.12. Curiously, France did not take to SMS in the same way, sending just under 20 messages on average per user per month. France has the same GSM technology as other European countries, so the uptake is not hampered by technical restrictions.
In the Republic of Ireland, in 2012, 1.5 billion messages were sent every quarter, on average 114 messages per person per month. In the United Kingdom, over 1 billion text messages were sent every week. The Eurovision Song Contest organized the first pan-European SMS voting in 2002, as a part of the voting system (there was also a voting over traditional landline phone lines). In 2005, the Eurovision Song Contest organized the biggest televoting ever (with SMS and phone voting). During roaming, (that is, when a user connects to another network in different country from their own) the prices may be higher, but in July 2009, EU legislation went into effect limiting this price to €0.11.
Mobile service providers in Finland offered contracts in which users can send 1000 text messages a month for €10. In Finland, which has very high mobile phone ownership rates, some TV channels began "SMS chat", which involved sending short messages to a phone number, and the messages would be shown on TV. Chats are always moderated, which prevents users from sending offensive material to the channel. The craze evolved into quizzes and strategy games and then faster-paced games designed for television and SMS control. Games require users to register their nicknames and send short messages to control a character onscreen. Messages usually cost 0.05 to 0.86 Euro apiece, and games can require the player to send dozens of messages. In December 2003, a Finnish TV channel, MTV3, put a Santa Claus character on-air reading aloud text messages sent in by viewers. On 12 March 2004, the first entirely "interactive" TV channel, VIISI, began operation in Finland. However, SBS Finland Oy took over the channel and turned it into a music channel named The Voice in November 2004. In 2006, the Prime Minister of Finland, Matti Vanhanen, made the news when he allegedly broke up with his girlfriend with a text message. In 2007, the first book written solely in text messages, Viimeiset viestit (Last Messages), was released by Finnish author Hannu Luntiala. It is about an executive who travels through Europe and India.
United States
In the United States, text messaging is very popular; as reported by CTIA in December 2009, the 286 million US subscribers sent 152.7 billion text messages per month, for an average of 534 messages per subscriber per month. The Pew Research Center found in May 2010 that 72% of U.S. adult cellphone users send and receive text messages. CTIA reported in 2022 that 2 trillion SMS and MMS were sent in the United States in 2021, showing continued popularity of the technology.
In the U.S., SMS is often charged both at the sender and at the destination, but, unlike phone calls, it cannot be rejected or dismissed. The reasons for lower uptake than other countries are varied. Many users have unlimited "mobile-to-mobile" minutes, high monthly minute allotments, or unlimited service. Moreover, "push to talk" services offer the instant connectivity of SMS and are typically unlimited. The integration between competing providers and technologies necessary for cross-network text messaging was not initially available. Some providers originally charged extra for texting, reducing its appeal. In the third quarter of 2006, at least 12 billion text messages were sent on AT&T's network, up almost 15% from the preceding quarter.
While texting is mainly popular among people from 13 to 22 years old, it is also increasing among adults and business users. The age that a child receives their first cell phone has also decreased, making text messaging a popular way of communicating. The number of texts sent in the US has gone up over the years as the price has gone down to an average of $0.10 per text sent and received. To convince more customers to buy unlimited text messaging plans, some major cellphone providers have increased the price to send and receive text messages from $.15 to $.20 per message. This is over $1,300 per megabyte. Many providers offer unlimited plans, which can result in a lower rate per text, given sufficient volume.
Japan
Japan was among the first countries to adopt short messages widely, with pioneering non-GSM services including J-Phone's SkyMail and NTT Docomo's Short Mail. Japanese adolescents first began text messaging, because it was a cheaper form of communication than the other available forms. Thus, Japanese theorists created the selective interpersonal relationship theory, claiming that mobile phones can change social networks among young people (classified as 13- to 30-year-olds). They theorized this age group had extensive but low-quality relationships with friends, and mobile phone usage may facilitate improvement in the quality of their relationships. They concluded this age group prefers "selective interpersonal relationships in which they maintain particular, partial, but rich relations, depending on the situation". The same studies showed participants rated friendships in which they communicated face-to-face and through text messaging as being more intimate than those in which they communicated solely face-to-face. This indicates participants make new relationships with face-to-face communication at an early stage, but use text messaging to increase their contact later on. As the relationships between participants grew more intimate, the frequency of text messaging also increased. However, short messaging has been largely rendered obsolete by the prevalence of mobile Internet e-mail, which can be sent to and received from any e-mail address, mobile or otherwise. That said, while usually presented to the user simply as a uniform "mail" service (and most users are unaware of the distinction), the operators may still internally transmit the content as short messages, especially if the destination is on the same network.
China
Text messaging has historically been popular and cheap in China. About 700 billion messages were sent in 2007. Text message spam has also been a problem in China. In 2007, 353.8 billion spam messages were sent, up 93% from the previous year. It is about 12.44 messages per week per person. In 2010, it was routine that the People's Republic of China government monitored text messages across the country for illegal content.
Among Chinese migrant workers with little formal education, it is common to refer to SMS manuals when text messaging. These manuals are published as cheap, smaller-than-pocket-size booklets that offer diverse linguistic phrases to utilize as messages.
Philippines
SMS was introduced to selected markets in the Philippines in 1995. In 1998, Philippine mobile service providers launched SMS more widely across the country, with initial television marketing campaigns targeting hearing-impaired users. The service was initially free with subscriptions, but Filipinos quickly exploited the feature to communicate for free instead of using voice calls, which they would be charged for. After telephone companies realized this trend, they began charging for SMS. The rate across networks is 1 peso per SMS (about US$0.023). Even after users were charged for SMS, it remained cheap, about one-tenth of the price of a voice call. This low price led to about five million Filipinos owning a cell phone by 2001. Because of the highly social nature of Philippine culture and the affordability of SMS compared to voice calls, SMS usage shot up. Filipinos used texting not only for social messages but also for political purposes, as it allowed the Filipinos to express their opinions on current events and political issues. It became a powerful tool for Filipinos in promoting or denouncing issues and was a key factor during the 2001 EDSA II revolution, which overthrew then-President Joseph Estrada, who was eventually found guilty of corruption. According to 2009 statistics, there were about 72 million mobile service subscriptions (roughly 80% of the Filipino population), with around 1.39 billion SMS messages being sent daily. Because of the large number of text messages being sent, the Philippines became known as the "text capital of the world" during the late 1990s until the early 2000s.
New Zealand
There are three mobile network companies operating in New Zealand, with some sub-brands and MVNOs.
Spark NZ (formerly Telecom NZ), was the first telecommunication company in New Zealand. In 2011, Spark was broken into two companies by regulation, with Chorus Ltd taking the landline infrastructure and Spark NZ providing services including over their mobile network. Vodafone NZ (now One NZ) acquired mobile network provider Bellsouth New Zealand in 1998 and had 2.32 million customers as of July 2013. Vodafone launched the first Text messaging service in 1999 and has introduced innovative TXT services like SafeTXT and CallMe 2degrees Mobile Ltd launched in August 2009.
In 2005, around 85% of the adult population had a mobile phone. In general, texting is more popular than making phone calls, as it is viewed as less intrusive and therefore more polite.
Sub-Saharan Africa
In 2009, it was predicted that text messaging would become a key revenue driver for mobile network operators in Africa over the following couple of years. Today, text messaging is already slowly gaining influence in the African market. One such person used text messaging to spread the word about HIV and AIDS. In September 2009, a multi-country campaign in Africa used text messaging to expose stock-outs of essential medicines at public health facilities and put pressure on governments to address the issue.
Social effects
The advent of text messaging made possible new forms of interaction that were not possible before. A person could carry out a conversation with another user without the constraint of being expected to reply within a short amount of time and without needing to set time aside to engage in conversation. With voice calling, both participants need to be free at the same time. Mobile phone users can maintain communication during situations in which a voice call is impractical, impossible, or unacceptable, such as during a school class or work meeting. Texting has provided a venue for participatory culture, allowing viewers to vote in online and TV polls, as well as receive information while they are on the move. Texting can also bring people together and create a sense of community through "Smart Mobs" or "Net War", which create "people power".
Research in 2015 has also proven that text messaging is somehow making the social distances larger and could be ruining verbal communication skills for many people.
Effect on language
The small phone keypad and the rapidity of typical text message exchanges have caused a number of spelling abbreviations: as in the phrase "txt msg", "u" (an abbreviation for "you"), "HMU"("hit me up"; i.e., call me), or use of camel case, such as in "ThisIsVeryLame". To avoid the even more limited message lengths allowed when using Cyrillic or Greek letters, speakers of languages written in those alphabets often use the Latin alphabet for their own language. In certain languages utilizing diacritic marks, such as Polish, SMS technology created an entire new variant of written language: characters normally written with diacritic marks (e.g., ą, ę, ś, ż in Polish) are now being written without them (as a, e, s, z) to enable using cell phones without Polish script or to save space in Unicode messages. Historically, this language developed out of shorthand used in bulletin board systems and later in Internet chat rooms, where users would abbreviate some words to allow a response to be typed more quickly, though the amount of time saved was often inconsequential. However, this became much more pronounced in SMS, where mobile phone users either have a numeric keyboard (with older cellphones) or a small QWERTY keyboard (for 2010s-era smartphones), so more effort is required to type each character, and there is sometimes a limit on the number of characters that may be sent. In Mandarin Chinese, numbers that sound similar to words are used in place of those words. For example, the numbers 520 in Chinese (wǔ èr líng) sound like the words for "I love you" (wǒ ài nǐ). The sequence 748 (qī sì bā) sounds like the curse "go to hell" (qù sǐ ba).
Predictive text software, which attempts to guess words (Tegic's T9 as well as iTap) or letters (Eatoni's LetterWise) reduces the labour of time-consuming input. This makes abbreviations not only less necessary but slower to type than regular words that are in the software's dictionary. However, it makes the messages longer, often requiring the text message to be sent in multiple parts and, therefore, costing more to send. The use of text messaging has changed the way that people talk and write essays, some believing it to be harmful. Children today are receiving cell phones at an age as young as eight years old; more than 35 per cent of children in second and third grade have their own mobile phones. Because of this, the texting language is integrated into the way that students think from an earlier age than ever before. In November 2006, New Zealand Qualifications Authority approved the move that allowed students of secondary schools to use mobile phone text language in the end-of-the-year-exam papers. Highly publicized reports, beginning in 2002, of the use of text language in school assignments, caused some to become concerned that the quality of written communication is on the decline, and other reports claim that teachers and professors are beginning to have a hard time controlling the problem. However, the notion that text language is widespread or harmful is refuted by research from linguistic experts.
An article in The New Yorker explores how text messaging has anglicized some of the world's languages. The use of diacritic marks is dropped in languages such as French, as well as symbols in Ethiopian languages. In his book, Txtng: the Gr8 Db8 (which translates as "Texting: the Great Debate"), David Crystal states that texters in all eleven languages use "lol" ("laughing out loud"), "u", "brb" ("be right back"), and "gr8" ("great"), all English-based shorthands. The use of pictograms and logograms in texts are present in every language. They shorten words by using symbols to represent the word or symbols whose name sounds like a syllable of the word such as in 2day or b4. This is commonly used in other languages as well. Crystal gives some examples in several languages such as Italian sei, "six", is used for sei, "you are". Example: dv6 = dove sei ("where are you") and French k7 = cassette ("cassette tape"). There is also the use of numeral sequences, substituting for several syllables of a word and creating whole phrases using numerals. For example, in French, a12c4 can be said as à un de ces quatres, "see you around" (literally: "to one of these four [days]"). An example of using symbols in texting and borrowing from English is the use of @. Whenever it is used in texting, its intended use is with the English pronunciation. Crystal gives the example of the Welsh use of @ in @F, pronounced ataf, meaning "to me". In character-based languages such as Chinese and Japanese, numbers are assigned syllables based on the shortened form of the pronunciation of the number, sometimes the English pronunciation of the number. In this way, numbers alone can be used to communicate whole passages, such as in Chinese, "8807701314520" () can be literally translated as "Hug hug you, kiss you, whole life, whole life I love you." English influences worldwide texting in variation, but still in combination with the individual properties of languages.
American popular culture is also recognized in shorthand. For example, Homer Simpson translates into: ~(_8^(|). Crystal also suggests that texting has led to more creativity in the English language, giving people opportunities to create their own slang, emoticons, abbreviations, acronyms, etc. The feeling of individualism and freedom makes texting more popular and a more efficient way to communicate. Crystal has also been quoted in saying that "In a logical world, text messaging should not have survived." But text messaging didn't just come out of nowhere. It originally began as a messaging system that would send out emergency information. But it gained immediate popularity with the public. What followed is the SMS we see today, which is a very quick and efficient way of sharing information from person to person. Work by Richard Ling has shown that texting has a gendered dimension and it plays into the development of teen identity. In addition we text to a very small number of other persons. For most people, half of their texts go to 3 – 5 other people.
Research by Rosen et al. (2009) found that those young adults who used more language-based textisms (shortcuts such as LOL, 2nite, etc.) in daily writing produced worse formal writing than those young adults who used fewer linguistic textisms in daily writing. However, the exact opposite was true for informal writing. This suggests that perhaps the act of using textisms to shorten communication words leads young adults to produce more informal writing, which may then help them to be better "informal" writers. Due to text messaging, teens are writing more, and some teachers see that this comfort with language can be harnessed to make better writers. This new form of communication may be encouraging students to put their thoughts and feelings into words and this may be able to be used as a bridge, to get them more interested in formal writing.
Joan H. Lee in her thesis, What does txting do 2 language: The influences of exposure to messaging and print media on acceptability constraints (2011), she associates exposure to text messaging with more rigid acceptability constraints. The thesis suggests that more exposure to the colloquial, Generation Text language of text messaging contributes to being less accepting of words. In contrast, Lee found that students with more exposure to traditional print media (such as books and magazines) were more accepting of both real and fictitious words. The thesis, which garnered international media attention, also presents a literature review of academic literature on the effects of text messaging on language. Texting has also been shown to have had no effect or some positive effects on literacy. According to Plester, Wood and Joshi and their research done on the study of 88 British 10–12-year-old children and their knowledge of text messages, "textisms are essentially forms of phonetic abbreviation" that show that "to produce and read such abbreviations arguably requires a level of phonological awareness (and orthographic awareness) in the child concerned".
Texting while driving
Texting while driving leads to increased distraction behind the wheel and can lead to an increased risk of an accident. In 2006, Liberty Mutual Insurance Group conducted a survey with more than 900 teens from over 26 high schools nationwide. The results showed that 87% of students found texting to be "very" or "extremely" distracting. A study by AAA found that 46% of teens admitted to being distracted behind the wheel due to texting. One example of distraction behind the wheel is the 2008 Chatsworth train collision, which killed 25 passengers. The engineer had sent 45 text messages while operating the train. A 2009 experiment with Car and Driver editor Eddie Alterman (that took place at a deserted airfield, for safety reasons) compared texting with drunk driving. The experiment found that texting while driving was more dangerous than being drunk. While being legally drunk added 4 feet to Alterman's stopping distance while going , reading an e-mail on a phone added , and sending a text message added .
In 2009, the Virginia Tech Transportation Institute released the results of an 18-month study that involved placing cameras inside the cabs of more than 100 long-haul trucks, which recorded the drivers over a combined driving distance of three million miles. The study concluded that when the drivers were texting, their risk of crashing was 23 times greater than when not texting.
Texting while walking
Due to the proliferation of smart phone applications performed while walking, "texting while walking" or "wexting" is the increasing practice of people being transfixed to their mobile device without looking in any direction but their personal screen while walking. First coined reference in 2015 in New York from Rentrak's chief client officer when discussing time spent with media and various media usage metrics. Text messaging among pedestrians leads to increased cognitive distraction and reduced situation awareness, and may lead to increases in unsafe behaviour leading to injury and death. Recent studies conducted on cell phone use while walking showed that cell phone users recall fewer objects when conversing, walk slower, have altered gait and are more unsafe when crossing a street. Additionally, some gait analyses showed that stance phase during overstepping motion, longitudinal and lateral deviation increased during cell phone operation, but step length and clearance did not; a different analysis did find increased step clearance and reduced step length.
It is unclear which processes may be affected by distraction, which types of distraction may affect which cognitive processes, and how individual differences may affect the influence of distraction. Lamberg and Muratori believe that engaging in a dual-task, such as texting while walking, may interfere with working memory and result in walking errors. Their study demonstrated that participants engaged in text messaging were unable to maintain walking speed or retain accurate spatial information, suggesting an inability to adequately divide their attention between two tasks. According to them, the addition of texting while walking with vision occluded increases the demands placed on the working memory system resulting in gait disruptions.
Texting on a phone distracts participants, even when the texting task used is a relatively simple one. Stavrinos et al. investigated the effect of other cognitive tasks, such as engaging in conversations or cognitive tasks on a phone, and found that participants actually have reduced visual awareness. This finding was supported by Licence et al., who conducted a similar study. For example, texting pedestrians may fail to notice unusual events in their environment, such as a unicycling clown. These findings suggest that tasks that require the allocation of cognitive resources can affect visual attention even when the task itself does not require the participants to avert their eyes from their environment. The act of texting itself seems to impair pedestrians' visual awareness. It appears that the distraction produced by texting is a combination of both a cognitive and visual perceptual distraction. A study conducted by Licence et al. supported some of these findings, particularly that those who text while walking significantly alter their gait. However, they also found that the gait pattern texters adopted was slower and more "protective", and consequently did not increase obstacle contact or tripping in a typical pedestrian context.
There have also been technological approaches to increase the safety/awareness of pedestrians that are (unintentionally) blind while using a smartphone, e.g., using a Kinect or an ultrasound phone cover as a virtual white cane, or using the built-in camera to algorithmically analyze single, respectively a stream of pictures for obstacles, with Wang et al. proposing to use machine learning to specifically detect incoming vehicles.
Sexting
Sexting is slang for the act of sending sexually explicit or suggestive content between mobile devices using SMS. It contains either text, images, or video that is intended to be sexually arousing. Sexting was reported as early as 2005 in The Sunday Telegraph Magazine, constituting a trend in the creative use of SMS to excite another with alluring messages throughout the day.
Although sexting often takes place consensually between two people, it can also occur against the wishes of a person who is the subject of the content. A number of instances have been reported in which the recipients of sexting have shared the content of the messages with others, with less intimate intentions, such as to impress their friends or embarrass their sender. Celebrities such as Miley Cyrus, Vanessa Hudgens, and Adrienne Bailon have been victims of such abuses of sexting.
A 2008 survey by The National Campaign to Prevent Teen and Unplanned Pregnancy and CosmoGirl.com suggested a trend of sexting and other seductive online content being readily shared between teens. One in five teen girls surveyed (22 per cent)—and 11 per cent of teen girls aged 13–16 years old—say they have electronically sent, or posted online, nude or semi-nude images of themselves. One-third (33 per cent) of teen boys and one-quarter (25 per cent) of teen girls say they were shown private nude or semi-nude images. According to the survey, sexually suggestive messages (text, e-mail, and instant messaging) were even more common than images, with 39 per cent of teens having sent or posted such messages, and half of the teens (50 per cent) having received them.
A 2012 study that has received wide international media attention was conducted at the University of Utah Department of Psychology by Donald S. Strassberg, Ryan Kelly McKinnon, Michael Sustaíta and Jordan Rullo. They surveyed 606 teenagers ages 14–18 and found that nearly 20 per cent of the students said they had sent a sexually explicit image of themselves via cell phone, and nearly twice as many said that they had received a sexually explicit picture. Of those receiving such a picture, over 25 per cent indicated that they had forwarded it to others. In addition, of those who had sent a sexually explicit picture, over a third had done so despite believing that there could be serious legal and other consequences if they got caught. Students who had sent a picture by cell phone were more likely than others to find the activity acceptable. The authors conclude: "These results argue for educational efforts such as cell phone safety assemblies, awareness days, integration into class curriculum and teacher training, designed to raise awareness about the potential consequences of sexting among young people." Sexting becomes a legal issue when teens (under 18) are involved, because any nude photos they may send of themselves would put the recipients in possession of child pornography.
In schools
Text messaging has affected students academically by creating an easier way to cheat on exams. In December 2002, a dozen students were caught cheating on an accounting exam through the use of text messages on their mobile phones. In December 2002, Hitotsubashi University in Japan failed 26 students for receiving emailed exam answers on their mobile phones. The number of students caught using mobile phones to cheat on exams has increased significantly in recent years. According to Okada (2005), most Japanese mobile phones can send and receive long text messages of between 250 and 3000 characters with graphics, video, audio, and Web links.
In England, 287 school and college students were excluded from exams in 2004 for using mobile phones during exams. Some teachers and professors claim that advanced texting features can lead to students cheating on exams. Students in high school and college classrooms are using their mobile phones to send and receive texts during lectures at high rates. Further, published research has established that students who text during college lectures have impaired memories of the lecture material compared to students who do not. For example, in one study, the number of irrelevant text messages sent and received during a lecture covering the topic of developmental psychology was related to students' memory of the lecture.
Bullying
Spreading rumors and gossip by text message, using text messages to bully individuals, or forwarding texts that contain defamatory content is an issue of great concern for parents and schools. Text "bullying" of this sort can cause distress and damage reputations. In some cases, individuals who are bullied online have committed suicide. Harding and Rosenberg (2005) argue that the urge to forward text messages can be difficult to resist, describing text messages as "loaded weapons".
Apple's messaging app, Messages, uses Apple's Internet-based messaging service, iMessage, to send messages to other iMessage users, and uses SMS as a fallback when no data connection is present, or when messaging non-iMessage users. It sets the color of messages depending on which technology was used. This has led to instances of iMessage users bullying people without iPhones.
Influence on perceptions of the student
When a student sends an email that contains phonetic abbreviations and acronyms that are common in text messaging (e.g., "gr8" instead of "great"), it can influence how that student is subsequently evaluated. In a study by Lewandowski and Harrington (2006), participants read a student's email sent to a professor that either contained text-messaging abbreviations (gr8, How R U?) or parallel text in standard English (great, How are you?), and then provided impressions of the sender. Students who used abbreviations in their email were perceived as having a less favorable personality and as putting forth less effort on an essay they submitted along with the email. Specifically, abbreviation users were seen as less intelligent, responsible, motivated, studious, dependable, and hard-working. These findings suggest that the nature of a student's email communication can influence how others perceive the student and their work.
However, students have become aware of the reality that using these textisms and adaptations can negatively impact their professionalism. Drouin and Davis surveyed American undergraduates in 2009 and found that three quarters of participants believed the use of textisms were not appropriate in formal messaging and writing. A study performed by Grace et al. (2013) asked 150 undergraduate students to rate the appropriateness of using textisms in a given scenario on a scale of one to five – five being entirely appropriate and one being not at all. All but eleven of the students rated the use of textisms in exams and typed assignments as "not at all appropriate", showing that the students are aware of how they must adapt their written language and tone depending on the context. Grace et al. (2010) went further, observing hundreds of academic papers from previous undergraduate students' exams, only to find that out of 533,500 words, a mere 0.02% were textisms. They owe this to the fact that the more accumulated experience a student has, the more they are able to understand when the "appropriate" and "inappropriate" times to use such language is.
Law and crime
Text messaging has been a subject of interest for police forces around the world. One of the issues of concern to law enforcement agencies is the use of encrypted text messages. In 2003, a British company developed a program called Fortress SMS which used 128 bit AES encryption to protect SMS messages. Police have also retrieved deleted text messages to aid them in solving crimes. For example, Swedish police retrieved deleted texts from a cult member who claimed she committed a double murder based on forwarded texts she received. Police in Tilburg, Netherlands, started an SMS alert program, in which they would send a message to ask citizens to be vigilant when a burglar was on the loose or a child was missing in their neighbourhood. Several thieves have been caught and children have been found using the SMS Alerts. The service has been expanding to other cities. A Malaysian–Australian company has released a multi-layer SMS security program. Boston police are now turning to text messaging to help stop crime. The Boston Police Department asks citizens to send texts to make anonymous crime tips.
Under some interpretations of sharia law, husbands can divorce their wives by the pronouncement of talaq. In 2003, a court in Malaysia upheld such a divorce pronouncement which was transmitted via SMS.
The Massachusetts Supreme Judicial Court ruled in 2017 that under the state constitution, police require a warrant before obtaining access to text messages without consent.
Social unrest
Texting has been used on a number of occasions with the result of the gathering of large aggressive crowds. SMS messaging drew a crowd to Cronulla Beach in Sydney resulting in the 2005 Cronulla riots. Not only were text messages circulating in the Sydney area but in other states as well (Daily Telegraph). The volume of such text messages and e-mails also increased in the wake of the riot. The crowd of 5,000 at stages became violent, attacking certain ethnic groups. Sutherland Shire Mayor directly blamed heavily circulated SMS messages for the unrest. NSW police considered whether people could be charged over the texting. Retaliatory attacks also used SMS.
The Narre Warren Incident, when a group of 500 party goers attended a party at Narre Warren in Melbourne, Australia, and rioted in January 2008, also was a response of communication being spread by SMS and Myspace. Following the incident, the Police Commissioner wrote an open letter asking young people to be aware of the power of SMS and the Internet. In Hong Kong, government officials find that text messaging helps socially because they can send multiple texts to the community. Officials say it is an easy way of contacting the community or individuals for meetings or events. Texting was used to coordinate gatherings during the 2009 Iranian election protests.
Between 2009 and 2012 the U.S. secretly created and funded a Twitter-like service for Cubans called ZunZuneo, initially based on mobile phone text message service and later with an internet interface. The service was funded by the U.S. Agency for International Development through its Office of Transition Initiatives, who utilized contractors and front companies in the Cayman Islands, Spain and Ireland. A longer-term objective was to organize "smart mobs" that might "renegotiate the balance of power between the state and society." A database about the subscribers was created, including gender, age, and "political tendencies". At its peak ZunZuneo had 40,000 Cuban users, but the service closed as financially unsustainable when U.S. funding was stopped.
In politics
Text messaging has affected the political world. American campaigns find that text messaging is a much easier, cheaper way of getting to the voters than the door-to-door approach. In 2006 Mexico's then president-elect Felipe Calderón launched millions of text messages in the days immediately preceding his narrow win over Andrés Manuel López Obrador. In January 2001, Joseph Estrada was forced to resign from the post of president of the Philippines. The popular campaign against him was widely reported to have been coordinated with SMS chain letters. A massive texting campaign was credited with boosting youth turnout in Spain's 2004 parliamentary elections. In 2008, Detroit Mayor Kwame Kilpatrick and his Chief of Staff at the time became entangled in a sex scandal stemming from the exchange of over 14,000 text messages that eventually led to his forced resignation, the conviction of perjury, and other charges.
Text messaging has been used to turn down other political leaders. During the 2004 U.S. Democratic and Republican National Conventions, protesters used an SMS-based organizing tool called TXTmob to get to opponents. In the last day before the 2004 presidential elections in Romania, a message against Adrian Năstase was largely circulated, thus breaking the laws that prohibited campaigning that day. Text messaging has helped politics by promoting campaigns.
On 20 January 2001, President Joseph Estrada of the Philippines became the first head of state in history to lose power to a smart mob. More than one million Manila residents assembled at the site of the 1986 People Power peaceful demonstrations that have toppled the Marcos regime. These people have organized themselves and coordinated their actions through text messaging. They were able to bring down a government without having to use any weapons or violence. Through text messaging, their plans and ideas were communicated to others and successfully implemented. Also, this move encouraged the military to withdraw their support from the regime, and as a result, the Estrada government fell. People were able to converge and unite with the use of their cell phones. "The rapid assembly of the anti-Estrada crowd was a hallmark of early smart mob technology, and the millions of text messages exchanged by the demonstrators in 2001 was, by all accounts, a key to the crowds esprit de corps."
Use in healthcare
Text messaging is a rapidly growing trend in Healthcare. A randomized controlled trial of text messaging intervention for diabetes in Bangladesh was one of the first robust trials to report improvement in diabetes management in a low-and-middle income country. A recent systematic review and individual participants data meta analysis from 3,779 participants reported that mobile phone text messaging could improve blood pressure and body mass index. Another study in people with type 2 diabetes showed that participants were willing to pay a modest amount to receive a diabetes text messaging program in addition to standard care. "One survey found that 73% of physicians text other physicians about work- similar to the overall percentage of the population that texts." A 2006 study of reminder messages sent to children and adolescents with type 1 diabetes mellitus showed favorable changes in adherence to treatment. A risk is that these physicians could be violating the Health Insurance Portability and Accountability Act. Where messages could be saved to a phone indefinitely, patient information could be subject to theft or loss, and could be seen by other unauthorized persons. The HIPAA privacy rule requires that any text message involving a medical decision must be available for the patient to access, meaning that any texts that are not documented in an EMR system could be a HIPAA violation.
Medical concerns
The excessive use of the thumb for pressing keys on mobile devices has led to a high rate of a form of repetitive strain injury termed "BlackBerry thumb" (although this refers to strain developed on older Blackberry devices, which had a scroll wheel on the side of the phone). An inflammation of the tendons in the thumb caused by constant text-messaging is also called text-messager's thumb, or texting tenosynovitis. Texting has also been linked as a secondary source in numerous traffic collisions, in which police investigations of mobile phone records have found that many drivers have lost control of their cars while attempting to send or retrieve a text message. Increasing cases of Internet addiction are now also being linked to text messaging, as mobile phones are now more likely to have e-mail and Web capabilities to complement the ability to text.
Etiquette
Texting etiquette refers to what is considered appropriate texting behaviour. These expectations may concern different areas, such as the context in which a text was sent and received/read, who each participant was with when the participant sent or received/read a text message or what constitutes impolite text messages. At the website of The Emily Post Institute, the topic of texting has spurred several articles with the "do's and dont's" regarding the new form of communication. One example from the site is: "Keep your message brief. No one wants to have an entire conversation with you by texting when you could just call him or her instead." Another example is: "Don't use all Caps. Typing a text message in all capital letters will appear as though you are shouting at the recipient, and should be avoided."
Expectations for etiquette may differ depending on various factors. For example, expectations for appropriate behaviour have been found to differ markedly between the U.S. and India. Another example is generational differences. In The M-Factor: How the Millennial Generation Is Rocking the Workplace, Lynne Lancaster and David Stillman note that younger Americans often do not consider it rude to answer their cell or begin texting in the middle of a face-to-face conversation with someone else, while older people, less used to the behavior and the accompanying lack of eye contact or attention, find this to be disruptive and ill-mannered. With regard to texting in the workplace, Plantronics studied how we communicate at work] and found that 58% of US knowledge workers have increased the use of text messaging for work in the past five years. The same study found that 33% of knowledge workers felt text messaging was critical or very important to success and productivity at work.
Typing awareness indicators
In some text messaging software products, an ellipsis is displayed while the interlocutor is typing characters. The feature has been referred to as a "typing awareness indicator", for which patents have been filed since the 1990s.
Challenges
Spam
In 2002, an increasing trend towards spamming mobile phone users through SMS prompted cellular-service carriers to take steps against the practice, before it became a widespread problem. No major spamming incidents involving SMS had been reported , but the existence of mobile phone spam has been noted by industry watchdogs including Consumer Reports magazine and the Utility Consumers' Action Network (UCAN). In 2005, UCAN brought a case against Sprint for spamming its customers and charging $0.10 per text message. The case was settled in 2006 with Sprint agreeing not to send customers Sprint advertisements via SMS. SMS expert Acision (formerly LogicaCMG Telecoms) reported a new type of SMS malice at the end of 2006, noting the first instances of SMiShing (a cousin to e-mail phishing scams). In SMiShing, users receive SMS messages posing to be from a company, enticing users to phone premium-rate numbers or reply with personal information. Similar concerns were reported by PhonepayPlus, a consumer watchdog in the United Kingdom, in 2012.
Pricing concerns
Concerns have been voiced over the excessive cost of off-plan text messaging in the United States. AT&T Mobility, along with most other service providers, charges texters 20 cents per message if they do not have a messaging plan or if they have exceeded their allotted number of texts. Given that an SMS message is at most 160 bytes in size, this cost scales to a cost of $1,310 per megabyte sent via text message. This is in sharp contrast with the price of unlimited data plans offered by the same carriers, which allow the transmission of hundreds of megabytes of data for monthly prices of about $15 to $45 in addition to a voice plan. As a comparison, a one-minute phone call uses up the same amount of network capacity as 600 text messages, meaning that if the same cost-per-traffic formula were applied to phone calls, cell phone calls would cost $120 per minute. With service providers gaining more customers and expanding their capacity, their overhead costs should be decreasing, not increasing. In 2005, text messaging generated nearly 70 billion dollars in revenue, as reported by Gartner, industry analysts, three times as much as Hollywood box office sales in 2005. World figures showed that over a trillion text messages were sent in 2005.
Although major cellphone providers deny any collusion, fees for out-of-package text messages have increased, doubling from 10 to 20 cents in the United States between 2007 and 2008 alone. On 16 July 2009, Senate hearings were held to look into any breach of the Sherman Antitrust Act. The same trend is visible in other countries, though increasingly widespread flat-rate plans, for example in Germany, do make text messaging easier, text messages sent abroad still result in higher costs.
Increasing competition
While text messaging is still a growing market, traditional SMS is becoming increasingly challenged by alternative messaging services which are available on smartphones with data connections. These services are much cheaper and offer more functionality like exchanging multimedia content (e.g. photos, videos or audio notes) and group messaging. Especially in western countries some of these services attract more and more users. Prominent examples of these include Apple's iMessage (exclusive to the Apple ecosystem) and the GSMA's RCS.
In 2021, 8.4 trillion SMS messages were sent globally, compared to 18.25 trillion for WhatsApp alone.
Security concerns
Experts have advised business users not to use consumer SMS for confidential communication. The contents of common SMS messages are known to the network operator's systems and personnel. Therefore, consumer SMS is not an appropriate technology for secure communications. To address this issue, many companies use an SMS gateway provider based on SS7 connectivity to route the messages. The advantage of this international termination model is the ability to route data directly through SS7, which gives the provider visibility of the complete path of the SMS. This means SMS messages can be sent directly to and from recipients without having to go through the SMS-C of other mobile operators. This approach reduces the number of mobile operators that handle the message; however, experts have advised not to consider it as an end-to-end secure communication, as the content of the message is exposed to the SMS gateway provider.
An alternative approach is to use end-to-end security software that runs on both the sending and receiving device, where the original text message is transmitted in encrypted form as a consumer SMS. By using key rotation, the encrypted text messages stored under data retention laws at the network operator cannot be decrypted even if one of the devices is compromised. A problem with this approach is that communicating devices needs to run compatible software. Failure rates without backward notification can be high between carriers. International texting can be unreliable depending on the country of origin, destination and respective operators. Differences in the character sets used for coding can cause a text message sent from one country to another to become unreadable.
In popular culture
Records and competition
The Guinness Book of World Records has a world record for text messaging, currently held by Sonja Kristiansen of Norway. Kristiansen keyed in the official text message, as established by Guinness, in 37.28 seconds. The message is, "The razor-toothed piranhas of the genera Serrasalmus and Pygocentrus are the most ferocious freshwater fish in the world. In reality, they seldom attack a human." In 2005, the record was held by a 24-year-old Scottish man, Craig Crosbie, who completed the same message in 48 seconds, beating the previous time by 19 seconds. The Book of Alternative Records lists Chris Young of Salem, Oregon as the world-record holder for the fastest 160-character text message where the contents of the message are not provided ahead of time. His record of 62.3 seconds was set on 23 May 2007.
Elliot Nicholls of Dunedin, New Zealand currently holds the world record for the fastest blindfolded text messaging. A record of a 160-letter text in 45 seconds while blindfolded was set on 17 November 2007, beating the old record of 1-minute 26 seconds set by an Italian in September 2006. Andrew Acklin of Ohio is credited with the world record for most text messages sent or received in a single month, with 200,052. His accomplishments were first in the World Records Academy and later followed up by Ripley's Believe It Or Not 2010: Seeing Is Believing. He has been acknowledged by The Universal Records Database for the most text messages in a single month; however, this has since been broken twice and as of 2010 was listed as 566607 messages by Fred Lindgren.
In January 2010, LG Electronics sponsored an international competition, the LG Mobile World Cup, to determine the fastest pair of texters. The winners were a team from South Korea, Ha Mok-min and Bae Yeong-ho. On 6 April 2011, SKH Apps released an iPhone app, iTextFast, to allow consumers to test their texting speed and practice the paragraph used by Guinness Book of World Records. As of 2011, best time listed on Game Center for that paragraph is 34.65 seconds.
Morse code
A few competitions have been held between expert Morse code operators and expert SMS users. Several mobile phones have Morse code ring tones and alert messages. For example, many Nokia mobile phones have an option to beep "S M S" in Morse code when it receives a short message. Some of these phones could also play the Nokia slogan "Connecting people" in Morse code as a message tone. There are third-party applications available for some mobile phones that allow Morse input for short messages.
Tattle texting
"Tattle texting" can mean either of two different texting trends:
Arena security
Many sports arenas now offer a number where patrons can text report security concerns, like drunk or unruly fans, or safety issues like spills. These programs have been praised by patrons and security personnel as more effective than traditional methods. For instance, the patron doesn't need to leave his seat and miss the event in order to report something important. Also, disruptive fans can be reported with relative anonymity. "Text tattling" also gives security personnel a useful tool to prioritize messages. For instance, a single complaint in one section about an unruly fan can be addressed when convenient, while multiple complaints by several different patrons can be acted upon immediately.
Smart cars
In this context, "tattle texting" refers to an automatic text sent by the computer in an automobile, because a preset condition was met. The most common use for this is for parents to receive texts from the car their child is driving, alerting them to speeding or other issues. Employers can also use the service to monitor their corporate vehicles. The technology is still new and (currently) only available on a few car models.
Common conditions that can be chosen to send a text are:
Speeding. With the use of GPS, stored maps, and speed limit information, the onboard computer can determine if the driver is exceeding the current speed limit. The device can store this information and/or send it to another recipient.
Range. Parents/employers can set a maximum range from a fixed location after which a "tattle text" is sent. Not only can this keep children close to home and keep employees from using corporate vehicles inappropriately, but it can also be a crucial tool for quickly identifying stolen vehicles, car jackings, and kidnappings.
| Technology | Media and communication | null |
305924 | https://en.wikipedia.org/wiki/Computational%20fluid%20dynamics | Computational fluid dynamics | Computational fluid dynamics (CFD) is a branch of fluid mechanics that uses numerical analysis and data structures to analyze and solve problems that involve fluid flows. Computers are used to perform the calculations required to simulate the free-stream flow of the fluid, and the interaction of the fluid (liquids and gases) with surfaces defined by boundary conditions. With high-speed supercomputers, better solutions can be achieved, and are often required to solve the largest and most complex problems. Ongoing research yields software that improves the accuracy and speed of complex simulation scenarios such as transonic or turbulent flows. Initial validation of such software is typically performed using experimental apparatus such as wind tunnels. In addition, previously performed analytical or empirical analysis of a particular problem can be used for comparison. A final validation is often performed using full-scale testing, such as flight tests.
CFD is applied to a wide range of research and engineering problems in many fields of study and industries, including aerodynamics and aerospace analysis, hypersonics, weather simulation, natural science and environmental engineering, industrial system design and analysis, biological engineering, fluid flows and heat transfer, engine and combustion analysis, and visual effects for film and games.
Background and history
The fundamental basis of almost all CFD problems is the Navier–Stokes equations, which define many single-phase (gas or liquid, but not both) fluid flows. These equations can be simplified by removing terms describing viscous actions to yield the Euler equations. Further simplification, by removing terms describing vorticity yields the full potential equations. Finally, for small perturbations in subsonic and supersonic flows (not transonic or hypersonic) these equations can be linearized to yield the linearized potential equations.
Historically, methods were first developed to solve the linearized potential equations. Two-dimensional (2D) methods, using conformal transformations of the flow about a cylinder to the flow about an airfoil were developed in the 1930s.
One of the earliest type of calculations resembling modern CFD are those by Lewis Fry Richardson, in the sense that these calculations used finite differences and divided the physical space in cells. Although they failed dramatically, these calculations, together with Richardson's book Weather Prediction by Numerical Process, set the basis for modern CFD and numerical meteorology. In fact, early CFD calculations during the 1940s using ENIAC used methods close to those in Richardson's 1922 book.
The computer power available paced development of three-dimensional methods. Probably the first work using computers to model fluid flow, as governed by the Navier–Stokes equations, was performed at Los Alamos National Lab, in the T3 group. This group was led by Francis H. Harlow, who is widely considered one of the pioneers of CFD. From 1957 to late 1960s, this group developed a variety of numerical methods to simulate transient two-dimensional fluid flows, such as particle-in-cell method, fluid-in-cell method, vorticity stream function method, and
marker-and-cell method. Fromm's vorticity-stream-function method for 2D, transient, incompressible flow was the first treatment of strongly contorting incompressible flows in the world.
The first paper with three-dimensional model was published by John Hess and A.M.O. Smith of Douglas Aircraft in 1967. This method discretized the surface of the geometry with panels, giving rise to this class of programs being called Panel Methods. Their method itself was simplified, in that it did not include lifting flows and hence was mainly applied to ship hulls and aircraft fuselages. The first lifting Panel Code (A230) was described in a paper written by Paul Rubbert and Gary Saaris of Boeing Aircraft in 1968. In time, more advanced three-dimensional Panel Codes were developed at Boeing (PANAIR, A502), Lockheed (Quadpan), Douglas (HESS), McDonnell Aircraft (MACAERO), NASA (PMARC) and Analytical Methods (WBAERO, USAERO and VSAERO). Some (PANAIR, HESS and MACAERO) were higher order codes, using higher order distributions of surface singularities, while others (Quadpan, PMARC, USAERO and VSAERO) used single singularities on each surface panel. The advantage of the lower order codes was that they ran much faster on the computers of the time. Today, VSAERO has grown to be a multi-order code and is the most widely used program of this class. It has been used in the development of many submarines, surface ships, automobiles, helicopters, aircraft, and more recently wind turbines. Its sister code, USAERO is an unsteady panel method that has also been used for modeling such things as high speed trains and racing yachts. The NASA PMARC code from an early version of VSAERO and a derivative of PMARC, named CMARC, is also commercially available.
In the two-dimensional realm, a number of Panel Codes have been developed for airfoil analysis and design. The codes typically have a boundary layer analysis included, so that viscous effects can be modeled. developed the PROFILE code, partly with NASA funding, which became available in the early 1980s. This was soon followed by Mark Drela's XFOIL code. Both PROFILE and XFOIL incorporate two-dimensional panel codes, with coupled boundary layer codes for airfoil analysis work. PROFILE uses a conformal transformation method for inverse airfoil design, while XFOIL has both a conformal transformation and an inverse panel method for airfoil design.
An intermediate step between Panel Codes and Full Potential codes were codes that used the Transonic Small Disturbance equations. In particular, the three-dimensional WIBCO code, developed by Charlie Boppe of Grumman Aircraft in the early 1980s has seen heavy use.
Developers turned to Full Potential codes, as panel methods could not calculate the non-linear flow present at transonic speeds. The first description of a means of using the Full Potential equations was published by Earll Murman and Julian Cole of Boeing in 1970. Frances Bauer, Paul Garabedian and David Korn of the Courant Institute at New York University (NYU) wrote a series of two-dimensional Full Potential airfoil codes that were widely used, the most important being named Program H. A further growth of Program H was developed by Bob Melnik and his group at Grumman Aerospace as Grumfoil. Antony Jameson, originally at Grumman Aircraft and the Courant Institute of NYU, worked with David Caughey to develop the important three-dimensional Full Potential code FLO22 in 1975. Many Full Potential codes emerged after this, culminating in Boeing's Tranair (A633) code, which still sees heavy use.
The next step was the Euler equations, which promised to provide more accurate solutions of transonic flows. The methodology used by Jameson in his three-dimensional FLO57 code (1981) was used by others to produce such programs as Lockheed's TEAM program and IAI/Analytical Methods' MGAERO program. MGAERO is unique in being a structured cartesian mesh code, while most other such codes use structured body-fitted grids (with the exception of NASA's highly successful CART3D code, Lockheed's SPLITFLOW code and Georgia Tech's NASCART-GT). Antony Jameson also developed the three-dimensional AIRPLANE code which made use of unstructured tetrahedral grids.
In the two-dimensional realm, Mark Drela and Michael Giles, then graduate students at MIT, developed the ISES Euler program (actually a suite of programs) for airfoil design and analysis. This code first became available in 1986 and has been further developed to design, analyze and optimize single or multi-element airfoils, as the MSES program. MSES sees wide use throughout the world. A derivative of MSES, for the design and analysis of airfoils in a cascade, is MISES, developed by Harold Youngren while he was a graduate student at MIT.
The Navier–Stokes equations were the ultimate target of development. Two-dimensional codes, such as NASA Ames' ARC2D code first emerged. A number of three-dimensional codes were developed (ARC3D, OVERFLOW, CFL3D are three successful NASA contributions), leading to numerous commercial packages.
Recently CFD methods have gained traction for modeling the flow behavior of granular materials within various chemical processes in engineering. This approach has emerged as a cost-effective alternative, offering a nuanced understanding of complex flow phenomena while minimizing expenses associated with traditional experimental methods.
Hierarchy of fluid flow equations
CFD can be seen as a group of computational methodologies (discussed below) used to solve equations governing fluid flow. In the application of CFD, a critical step is to decide which set of physical assumptions and related equations need to be used for the problem at hand. To illustrate this step, the following summarizes the physical assumptions/simplifications taken in equations of a flow that is single-phase (see multiphase flow and two-phase flow), single-species (i.e., it consists of one chemical species), non-reacting, and (unless said otherwise) compressible. Thermal radiation is neglected, and body forces due to gravity are considered (unless said otherwise). In addition, for this type of flow, the next discussion highlights the hierarchy of flow equations solved with CFD. Note that some of the following equations could be derived in more than one way.
Conservation laws (CL): These are the most fundamental equations considered with CFD in the sense that, for example, all the following equations can be derived from them. For a single-phase, single-species, compressible flow one considers the conservation of mass, conservation of linear momentum, and conservation of energy.
Continuum conservation laws (CCL): Start with the CL. Assume that mass, momentum and energy are locally conserved: These quantities are conserved and cannot "teleport" from one place to another but can only move by a continuous flow (see continuity equation). Another interpretation is that one starts with the CL and assumes a continuum medium (see continuum mechanics). The resulting system of equations is unclosed since to solve it one needs further relationships/equations: (a) constitutive relationships for the viscous stress tensor; (b) constitutive relationships for the diffusive heat flux; (c) an equation of state (EOS), such as the ideal gas law; and, (d) a caloric equation of state relating temperature with quantities such as enthalpy or internal energy.
Compressible Navier-Stokes equations (C-NS): Start with the CCL. Assume a Newtonian viscous stress tensor (see Newtonian fluid) and a Fourier heat flux (see heat flux). The C-NS need to be augmented with an EOS and a caloric EOS to have a closed system of equations.
Incompressible Navier-Stokes equations (I-NS): Start with the C-NS. Assume that density is always and everywhere constant. Another way to obtain the I-NS is to assume that the Mach number is very small and that temperature differences in the fluid are very small as well. As a result, the mass-conservation and momentum-conservation equations are decoupled from the energy-conservation equation, so one only needs to solve for the first two equations.
Compressible Euler equations (EE): Start with the C-NS. Assume a frictionless flow with no diffusive heat flux.
Weakly compressible Navier-Stokes equations (WC-NS): Start with the C-NS. Assume that density variations depend only on temperature and not on pressure. For example, for an ideal gas, use , where is a conveniently-defined reference pressure that is always and everywhere constant, is density, is the specific gas constant, and is temperature. As a result, the WC-NS do not capture acoustic waves. It is also common in the WC-NS to neglect the pressure-work and viscous-heating terms in the energy-conservation equation. The WC-NS are also called the C-NS with the low-Mach-number approximation.
Boussinesq equations: Start with the C-NS. Assume that density variations are always and everywhere negligible except in the gravity term of the momentum-conservation equation (where density multiplies the gravitational acceleration). Also assume that various fluid properties such as viscosity, thermal conductivity, and heat capacity are always and everywhere constant. The Boussinesq equations are widely used in microscale meteorology.
Compressible Reynolds-averaged Navier–Stokes equations and compressible Favre-averaged Navier-Stokes equations (C-RANS and C-FANS): Start with the C-NS. Assume that any flow variable , such as density, velocity and pressure, can be represented as , where is the ensemble-average of any flow variable, and is a perturbation or fluctuation from this average. is not necessarily small. If is a classic ensemble-average (see Reynolds decomposition) one obtains the Reynolds-averaged Navier–Stokes equations. And if is a density-weighted ensemble-average one obtains the Favre-averaged Navier-Stokes equations. As a result, and depending on the Reynolds number, the range of scales of motion is greatly reduced, something which leads to much faster solutions in comparison to solving the C-NS. However, information is lost, and the resulting system of equations requires the closure of various unclosed terms, notably the Reynolds stress.
Ideal flow or potential flow equations: Start with the EE. Assume zero fluid-particle rotation (zero vorticity) and zero flow expansion (zero divergence). The resulting flowfield is entirely determined by the geometrical boundaries. Ideal flows can be useful in modern CFD to initialize simulations.
Linearized compressible Euler equations (LEE): Start with the EE. Assume that any flow variable , such as density, velocity and pressure, can be represented as , where is the value of the flow variable at some reference or base state, and is a perturbation or fluctuation from this state. Furthermore, assume that this perturbation is very small in comparison with some reference value. Finally, assume that satisfies "its own" equation, such as the EE. The LEE and its many variations are widely used in computational aeroacoustics.
Sound wave or acoustic wave equation: Start with the LEE. Neglect all gradients of and , and assume that the Mach number at the reference or base state is very small. The resulting equations for density, momentum and energy can be manipulated into a pressure equation, giving the well-known sound wave equation.
Shallow water equations (SW): Consider a flow near a wall where the wall-parallel length-scale of interest is much larger than the wall-normal length-scale of interest. Start with the EE. Assume that density is always and everywhere constant, neglect the velocity component perpendicular to the wall, and consider the velocity parallel to the wall to be spatially-constant.
Boundary layer equations (BL): Start with the C-NS (I-NS) for compressible (incompressible) boundary layers. Assume that there are thin regions next to walls where spatial gradients perpendicular to the wall are much larger than those parallel to the wall.
Bernoulli equation: Start with the EE. Assume that density variations depend only on pressure variations. See Bernoulli's Principle.
Steady Bernoulli equation: Start with the Bernoulli Equation and assume a steady flow. Or start with the EE and assume that the flow is steady and integrate the resulting equation along a streamline.
Stokes Flow or creeping flow equations: Start with the C-NS or I-NS. Neglect the inertia of the flow. Such an assumption can be justified when the Reynolds number is very low. As a result, the resulting set of equations is linear, something which simplifies greatly their solution.
Two-dimensional channel flow equation: Consider the flow between two infinite parallel plates. Start with the C-NS. Assume that the flow is steady, two-dimensional, and fully developed (i.e., the velocity profile does not change along the streamwise direction). Note that this widely-used fully-developed assumption can be inadequate in some instances, such as some compressible, microchannel flows, in which case it can be supplanted by a locally fully-developed assumption.
One-dimensional Euler equations or one-dimensional gas-dynamic equations (1D-EE): Start with the EE. Assume that all flow quantities depend only on one spatial dimension.
Fanno flow equation: Consider the flow inside a duct with constant area and adiabatic walls. Start with the 1D-EE. Assume a steady flow, no gravity effects, and introduce in the momentum-conservation equation an empirical term to recover the effect of wall friction (neglected in the EE). To close the Fanno flow equation, a model for this friction term is needed. Such a closure involves problem-dependent assumptions.
Rayleigh flow equation. Consider the flow inside a duct with constant area and either non-adiabatic walls without volumetric heat sources or adiabatic walls with volumetric heat sources. Start with the 1D-EE. Assume a steady flow, no gravity effects, and introduce in the energy-conservation equation an empirical term to recover the effect of wall heat transfer or the effect of the heat sources (neglected in the EE).
Methodology
In all of these approaches the same basic procedure is followed.
During preprocessing
The geometry and physical bounds of the problem can be defined using computer aided design (CAD). From there, data can be suitably processed (cleaned-up) and the fluid volume (or fluid domain) is extracted.
The volume occupied by the fluid is divided into discrete cells (the mesh). The mesh may be uniform or non-uniform, structured or unstructured, consisting of a combination of hexahedral, tetrahedral, prismatic, pyramidal or polyhedral elements.
The physical modeling is defined – for example, the equations of fluid motion + enthalpy + radiation + species conservation
Boundary conditions are defined. This involves specifying the fluid behaviour and properties at all bounding surfaces of the fluid domain. For transient problems, the initial conditions are also defined.
The simulation is started and the equations are solved iteratively as a steady-state or transient.
Finally a postprocessor is used for the analysis and visualization of the resulting solution.
Discretization methods
The stability of the selected discretisation is generally established numerically rather than analytically as with simple linear problems. Special care must also be taken to ensure that the discretisation handles discontinuous solutions gracefully. The Euler equations and Navier–Stokes equations both admit shocks and contact surfaces.
Some of the discretization methods being used are:
Finite volume method
The finite volume method (FVM) is a common approach used in CFD codes, as it has an advantage in memory usage and solution speed, especially for large problems, high Reynolds number turbulent flows, and source term dominated flows (like combustion).
In the finite volume method, the governing partial differential equations (typically the Navier-Stokes equations, the mass and energy conservation equations, and the turbulence equations) are recast in a conservative form, and then solved over discrete control volumes. This discretization guarantees the conservation of fluxes through a particular control volume. The finite volume equation yields governing equations in the form,
where is the vector of conserved variables, is the vector of fluxes (see Euler equations or Navier–Stokes equations), is the volume of the control volume element, and is the surface area of the control volume element.
Finite element method
The finite element method (FEM) is used in structural analysis of solids, but is also applicable to fluids. However, the FEM formulation requires special care to ensure a conservative solution. The FEM formulation has been adapted for use with fluid dynamics governing equations. Although FEM must be carefully formulated to be conservative, it is much more stable than the finite volume approach. FEM also provides more accurate solutions for smooth problems comparing to FVM. Another advantage of FEM is that it can handle complex geometries and boundary conditions. However, FEM can require more memory and has slower solution times than the FVM.
In this method, a weighted residual equation is formed:
where is the equation residual at an element vertex , is the conservation equation expressed on an element basis, is the weight factor, and is the volume of the element.
Finite difference method
The finite difference method (FDM) has historical importance and is simple to program. It is currently only used in few specialized codes, which handle complex geometry with high accuracy and efficiency by using embedded boundaries or overlapping grids (with the solution interpolated across each grid).
where is the vector of conserved variables, and , , and are the fluxes in the , , and directions respectively.
Spectral element method
Spectral element method is a finite element type method. It requires the mathematical problem (the partial differential equation) to be cast in a weak formulation. This is typically done by multiplying the differential equation by an arbitrary test function and integrating over the whole domain. Purely mathematically, the test functions are completely arbitrary - they belong to an infinite-dimensional function space. Clearly an infinite-dimensional function space cannot be represented on a discrete spectral element mesh; this is where the spectral element discretization begins. The most crucial thing is the choice of interpolating and testing functions. In a standard, low order FEM in 2D, for quadrilateral elements the most typical choice is the bilinear test or interpolating function of the form . In a spectral element method however, the interpolating and test functions are chosen to be polynomials of a very high order (typically e.g. of the 10th order in CFD applications). This guarantees the rapid convergence of the method. Furthermore, very efficient integration procedures must be used, since the number of integrations to be performed in numerical codes is big. Thus, high order Gauss integration quadratures are employed, since they achieve the highest accuracy with the smallest number of computations to be carried out.
At the time there are some academic CFD codes based on the spectral element method and some more are currently under development, since the new time-stepping schemes arise in the scientific world.
Lattice Boltzmann method
The lattice Boltzmann method (LBM) with its simplified kinetic picture on a lattice provides a computationally efficient description of hydrodynamics.
Unlike the traditional CFD methods, which solve the conservation equations of macroscopic properties (i.e., mass, momentum, and energy) numerically, LBM models the fluid consisting of fictive particles, and such particles perform consecutive propagation and collision processes over a discrete lattice mesh. In this method, one works with the discrete in space and time version of the kinetic evolution equation in the Boltzmann Bhatnagar-Gross-Krook (BGK) form.
Vortex method
The vortex method, also Lagrangian Vortex Particle Method, is a meshfree technique for the simulation of incompressible turbulent flows. In it, vorticity is discretized onto Lagrangian particles, these computational elements being called vortices, vortons, or vortex particles. Vortex methods were developed as a grid-free methodology that would not be limited by the fundamental smoothing effects associated with grid-based methods. To be practical, however, vortex methods require means for rapidly computing velocities from the vortex elements – in other words they require the solution to a particular form of the N-body problem (in which the motion of N objects is tied to their mutual influences). This breakthrough came in the 1980s with the development of the Barnes-Hut and fast multipole method (FMM) algorithms. These paved the way to practical computation of the velocities from the vortex elements.
Software based on the vortex method offer a new means for solving tough fluid dynamics problems with minimal user intervention. All that is required is specification of problem geometry and setting of boundary and initial conditions. Among the significant advantages of this modern technology;
It is practically grid-free, thus eliminating numerous iterations associated with RANS and LES.
All problems are treated identically. No modeling or calibration inputs are required.
Time-series simulations, which are crucial for correct analysis of acoustics, are possible.
The small scale and large scale are accurately simulated at the same time.
Boundary element method
In the boundary element method, the boundary occupied by the fluid is divided into a surface mesh.
High-resolution discretization schemes
High-resolution schemes are used where shocks or discontinuities are present. Capturing sharp changes in the solution requires the use of second or higher-order numerical schemes that do not introduce spurious oscillations. This usually necessitates the application of flux limiters to ensure that the solution is total variation diminishing.
Turbulence models
In computational modeling of turbulent flows, one common objective is to obtain a model that can predict quantities of interest, such as fluid velocity, for use in engineering designs of the system being modeled. For turbulent flows, the range of length scales and complexity of phenomena involved in turbulence make most modeling approaches prohibitively expensive; the resolution required to resolve all scales involved in turbulence is beyond what is computationally possible. The primary approach in such cases is to create numerical models to approximate unresolved phenomena. This section lists some commonly used computational models for turbulent flows.
Turbulence models can be classified based on computational expense, which corresponds to the range of scales that are modeled versus resolved (the more turbulent scales that are resolved, the finer the resolution of the simulation, and therefore the higher the computational cost). If a majority or all of the turbulent scales are not modeled, the computational cost is very low, but the tradeoff comes in the form of decreased accuracy.
In addition to the wide range of length and time scales and the associated computational cost, the governing equations of fluid dynamics contain a non-linear convection term and a non-linear and non-local pressure gradient term. These nonlinear equations must be solved numerically with the appropriate boundary and initial conditions.
Reynolds-averaged Navier–Stokes
Reynolds-averaged Navier–Stokes (RANS) equations are the oldest approach to turbulence modeling. An ensemble version of the governing equations is solved, which introduces new apparent stresses known as Reynolds stresses. This adds a second-order tensor of unknowns for which various models can provide different levels of closure. It is a common misconception that the RANS equations do not apply to flows with a time-varying mean flow because these equations are 'time-averaged'. In fact, statistically unsteady (or non-stationary) flows can equally be treated. This is sometimes referred to as URANS. There is nothing inherent in Reynolds averaging to preclude this, but the turbulence models used to close the equations are valid only as long as the time over which these changes in the mean occur is large compared to the time scales of the turbulent motion containing most of the energy.
RANS models can be divided into two broad approaches:
Boussinesq hypothesis This method involves using an algebraic equation for the Reynolds stresses which include determining the turbulent viscosity, and depending on the level of sophistication of the model, solving transport equations for determining the turbulent kinetic energy and dissipation. Models include k-ε (Launder and Spalding), Mixing Length Model (Prandtl), and Zero Equation Model (Cebeci and Smith). The models available in this approach are often referred to by the number of transport equations associated with the method. For example, the Mixing Length model is a "Zero Equation" model because no transport equations are solved; the is a "Two Equation" model because two transport equations (one for and one for ) are solved.
Reynolds stress model (RSM) This approach attempts to actually solve transport equations for the Reynolds stresses. This means introduction of several transport equations for all the Reynolds stresses and hence this approach is much more costly in CPU effort.
Large eddy simulation
Large eddy simulation (LES) is a technique in which the smallest scales of the flow are removed through a filtering operation, and their effect modeled using subgrid scale models. This allows the largest and most important scales of the turbulence to be resolved, while greatly reducing the computational cost incurred by the smallest scales. This method requires greater computational resources than RANS methods, but is far cheaper than DNS.
Detached eddy simulation
Detached eddy simulations (DES) is a modification of a RANS model in which the model switches to a subgrid scale formulation in regions fine enough for LES calculations. Regions near solid boundaries and where the turbulent length scale is less than the maximum grid dimension are assigned the RANS mode of solution. As the turbulent length scale exceeds the grid dimension, the regions are solved using the LES mode. Therefore, the grid resolution for DES is not as demanding as pure LES, thereby considerably cutting down the cost of the computation. Though DES was initially formulated for the Spalart-Allmaras model (Philippe R. Spalart et al., 1997), it can be implemented with other RANS models (Strelets, 2001), by appropriately modifying the length scale which is explicitly or implicitly involved in the RANS model. So while Spalart–Allmaras model based DES acts as LES with a wall model, DES based on other models (like two equation models) behave as a hybrid RANS-LES model. Grid generation is more complicated than for a simple RANS or LES case due to the RANS-LES switch. DES is a non-zonal approach and provides a single smooth velocity field across the RANS and the LES regions of the solutions.
Direct numerical simulation
Direct numerical simulation (DNS) resolves the entire range of turbulent length scales. This marginalizes the effect of models, but is extremely expensive. The computational cost is proportional to . DNS is intractable for flows with complex geometries or flow configurations.
Coherent vortex simulation
The coherent vortex simulation approach decomposes the turbulent flow field into a coherent part, consisting of organized vortical motion, and the incoherent part, which is the random background flow. This decomposition is done using wavelet filtering. The approach has much in common with LES, since it uses decomposition and resolves only the filtered portion, but different in that it does not use a linear, low-pass filter. Instead, the filtering operation is based on wavelets, and the filter can be adapted as the flow field evolves. Farge and Schneider tested the CVS method with two flow configurations and showed that the coherent portion of the flow exhibited the energy spectrum exhibited by the total flow, and corresponded to coherent structures (vortex tubes), while the incoherent parts of the flow composed homogeneous background noise, which exhibited no organized structures. Goldstein and Vasilyev applied the FDV model to large eddy simulation, but did not assume that the wavelet filter eliminated all coherent motions from the subfilter scales. By employing both LES and CVS filtering, they showed that the SFS dissipation was dominated by the SFS flow field's coherent portion.
PDF methods
Probability density function (PDF) methods for turbulence, first introduced by Lundgren, are based on tracking the one-point PDF of the velocity, , which gives the probability of the velocity at point being between and . This approach is analogous to the kinetic theory of gases, in which the macroscopic properties of a gas are described by a large number of particles. PDF methods are unique in that they can be applied in the framework of a number of different turbulence models; the main differences occur in the form of the PDF transport equation. For example, in the context of large eddy simulation, the PDF becomes the filtered PDF. PDF methods can also be used to describe chemical reactions, and are particularly useful for simulating chemically reacting flows because the chemical source term is closed and does not require a model. The PDF is commonly tracked by using Lagrangian particle methods; when combined with large eddy simulation, this leads to a Langevin equation for subfilter particle evolution.
Vorticity confinement method
The vorticity confinement (VC) method is an Eulerian technique used in the simulation of turbulent wakes. It uses a solitary-wave like approach to produce a stable solution with no numerical spreading. VC can capture the small-scale features to within as few as 2 grid cells. Within these features, a nonlinear difference equation is solved as opposed to the finite difference equation. VC is similar to shock capturing methods, where conservation laws are satisfied, so that the essential integral quantities are accurately computed.
Linear eddy model
The Linear eddy model is a technique used to simulate the convective mixing that takes place in turbulent flow. Specifically, it provides a mathematical way to describe the interactions of a scalar variable within the vector flow field. It is primarily used in one-dimensional representations of turbulent flow, since it can be applied across a wide range of length scales and Reynolds numbers. This model is generally used as a building block for more complicated flow representations, as it provides high resolution predictions that hold across a large range of flow conditions.
Two-phase flow
The modeling of two-phase flow is still under development. Different methods have been proposed, including the Volume of fluid method, the level-set method and front tracking. These methods often involve a tradeoff between maintaining a sharp interface or conserving mass . This is crucial since the evaluation of the density, viscosity and surface tension is based on the values averaged over the interface.
Solution algorithms
Discretization in the space produces a system of ordinary differential equations for unsteady problems and algebraic equations for steady problems. Implicit or semi-implicit methods are generally used to integrate the ordinary differential equations, producing a system of (usually) nonlinear algebraic equations. Applying a Newton or Picard iteration produces a system of linear equations which is nonsymmetric in the presence of advection and indefinite in the presence of incompressibility. Such systems, particularly in 3D, are frequently too large for direct solvers, so iterative methods are used, either stationary methods such as successive overrelaxation or Krylov subspace methods. Krylov methods such as GMRES, typically used with preconditioning, operate by minimizing the residual over successive subspaces generated by the preconditioned operator.
Multigrid has the advantage of asymptotically optimal performance on many problems. Traditional solvers and preconditioners are effective at reducing high-frequency components of the residual, but low-frequency components typically require many iterations to reduce. By operating on multiple scales, multigrid reduces all components of the residual by similar factors, leading to a mesh-independent number of iterations.
For indefinite systems, preconditioners such as incomplete LU factorization, additive Schwarz, and multigrid perform poorly or fail entirely, so the problem structure must be used for effective preconditioning. Methods commonly used in CFD are the SIMPLE and Uzawa algorithms which exhibit mesh-dependent convergence rates, but recent advances based on block LU factorization combined with multigrid for the resulting definite systems have led to preconditioners that deliver mesh-independent convergence rates.
Unsteady aerodynamics
CFD made a major break through in late 70s with the introduction of LTRAN2, a 2-D code to model oscillating airfoils based on transonic small perturbation theory by Ballhaus and associates. It uses a Murman-Cole switch algorithm for modeling the moving shock-waves. Later it was extended to 3-D with use of a rotated difference scheme by AFWAL/Boeing that resulted in LTRAN3.
Biomedical engineering
CFD investigations are used to clarify the characteristics of aortic flow in details that are beyond the capabilities of experimental measurements. To analyze these conditions, CAD models of the human vascular system are extracted employing modern imaging techniques such as MRI or Computed Tomography. A 3D model is reconstructed from this data and the fluid flow can be computed. Blood properties such as density and viscosity, and realistic boundary conditions (e.g. systemic pressure) have to be taken into consideration. Therefore, making it possible to analyze and optimize the flow in the cardiovascular system for different applications.
CPU versus GPU
Traditionally, CFD simulations are performed on CPUs.
In a more recent trend, simulations are also performed on GPUs. These typically contain slower but more processors. For CFD algorithms that feature good parallelism performance (i.e. good speed-up by adding more cores) this can greatly reduce simulation times. Fluid-implicit particle and lattice-Boltzmann methods are typical examples of codes that scale well on GPUs.
| Physical sciences | Fluid mechanics | Physics |
3694774 | https://en.wikipedia.org/wiki/Formal%20science | Formal science | Formal science is a branch of science studying disciplines concerned with abstract structures described by formal systems, such as logic, mathematics, statistics, theoretical computer science, artificial intelligence, information theory, game theory, systems theory, decision theory and theoretical linguistics. Whereas the natural sciences and social sciences seek to characterize physical systems and social systems, respectively, using empirical methods, the formal sciences use language tools concerned with characterizing abstract structures described by formal systems. The formal sciences aid the natural and social sciences by providing information about the structures used to describe the physical world, and what inferences may be made about them.
Branches
Logic (also a branch of philosophy)
Mathematics
Statistics
Systems science
Data science
Information science
Computer science
Cryptography
Differences from other sciences
Because of their non-empirical nature, formal sciences are construed by outlining a set of axioms and definitions from which other statements (theorems) are deduced. For this reason, in Rudolf Carnap's logical-positivist conception of the epistemology of science, theories belonging to formal sciences are understood to contain no synthetic statements, instead containing only analytic statements.
| Physical sciences | Science basics | Basics and measurement |
3696045 | https://en.wikipedia.org/wiki/African%20wild%20ass | African wild ass | The African wild ass (Equus africanus) or African wild donkey is a wild member of the horse family, Equidae. This species is thought to be the ancestor of the domestic donkey (Equus asinus), which is sometimes placed within the same species. They live in the deserts and other arid areas of the Horn of Africa, in Eritrea, Ethiopia and Somalia. It formerly had a wider range north and west into Sudan, Egypt, and Libya. It is Critically Endangered, with about 570 existing in the wild.
Description
The African wild ass is about tall and weighs approximately . The short, smooth coat is a light grey to fawn colour, fading quickly to white on the undersides and legs. There is a slender, dark dorsal stripe in all subspecies, while in the Nubian wild ass (E. a. africanus), as well as the domestic donkey, there is a stripe across the shoulder. The legs of the Somali wild ass (E. a. somaliensis) are horizontally striped with black, resembling those of a zebra. On the nape of the neck, there is a stiff, upright mane, the hairs of which are tipped with black. The ears are large with black margins. The tail terminates with a black brush. The hooves are slender and approximately of the diameter as the legs.
Evolution
The genus Equus, which includes all extant equines, is believed to have evolved from Dinohippus, via the intermediate form Plesippus. One of the oldest species is Equus simplicidens, described as zebra-like with a donkey-shaped head. The oldest fossil to date is ~3.5 million years old from Idaho in the United States. The genus appears to have spread quickly into the Old World, with the similarly aged Equus livenzovensis documented from western Europe and Russia.
Molecular phylogenies indicate the most recent common ancestor of all modern equids (members of the genus Equus) lived ~5.6 (3.9–7.8) mya. Direct paleogenomic sequencing of a 700,000-year-old middle Pleistocene horse metapodial bone from Canada implies a more recent 4.07 Myr before present date for the most recent common ancestor (MRCA) within the range of 4.0 to 4.5 Myr BP. The oldest divergencies are the Asian hemiones (subgenus E. (Asinus), including the kulan, onager, and kiang), followed by the African zebras (subgenera E. (Dolichohippus), and E. (Hippotigris)). All other modern forms including the domesticated horse (and many fossil Pliocene and Pleistocene forms) belong to the subgenus E. (Equus) which diverged ~4.8 (3.2–6.5) million years ago.
Taxonomy
Different authors consider the African wild ass and the domesticated donkey one or two species; either view is technically legitimate, though the former is phylogenetically more accurate. However, the American Society of Mammalogists classifies the donkey as a distinct species, as it does with almost all domestic mammals.
The species name for the African wild ass is sometimes given as asinus, from the domestic donkey, whose specific name is older and usually would have priority. But this usage is erroneous since the International Commission on Zoological Nomenclature has conserved the name Equus africanus in Opinion 2027. This was done to prevent the confusing situation of the phylogenetic ancestor being taxonomically included in its descendant.
Thus, if one species is recognized, the correct scientific name of the donkey is E. africanus asinus.
The first published name for the African wild ass, Asinus africanus, Fitzinger, 1858, is a nomen nudum. The name Equus taeniopus von Heuglin, 1861 is rejected as indeterminable, as it is based on an animal that cannot be identified and may have been a hybrid between a domestic donkey and a Somali wild ass; the type has not been preserved. The first available name thus becomes Asinus africanus von Heuglin & Fitzinger, 1866. A lectotype is designated: a skull of an adult female collected by von Heuglin near Atbarah River, Sudan, and present in the State Museum of Natural History Stuttgart, MNS 32026. The subspecies recognized:
Cladogram based on whole nuclear genomes after Özkan et al. 2024.African wild asses of unknown subspecies lived in the Sahara, Egypt, and Arabia, earlier in the Holocene.
Habitat
African wild asses are well suited to life in a desert or semidesert environment. They have tough digestive systems, which can break down desert vegetation and extract moisture from food efficiently. They can also go without water for a fairly long time. Their large ears give them an excellent sense of hearing and help in cooling.
Because of the sparse vegetation in their environment wild asses live somewhat separated from each other (except for mothers and young), unlike the tightly grouped herds of wild horses. They have very loud voices, which can be heard for over , which helps them to keep in contact with other asses over the wide spaces of the desert.
Behavior
The African wild ass is primarily active in the cooler hours between late afternoon and early morning, seeking shade and shelter amongst the rocky hills during the day. The Somali wild ass is also very agile and nimble-footed, capable of moving quickly across boulder fields and in the mountains. On the flat, it has been recorded reaching speeds of . In keeping with these feats, its soles are particularly hard and its hooves grow very quickly.
Mature males defend large territories around 23 square kilometres in size, marking them with dung heaps – an essential marker in the flat, monotonous terrain. Due to the size of these ranges, the dominant male cannot exclude other males. Rather, intruders are tolerated – recognized and treated as subordinates, and kept as far away as possible from any of the resident females. In the presence of estrous females, the males bray loudly. These animals live in loose herds of up to fifty individuals.
In the wild, African wild ass breeding occurs during the wet season. The gestation period lasts for 11 to 12 months, and one foal is born during the period from October to February. The foal weans for 6 to 8 months after birth, reaching sexual maturity at the age of 2 years. Lifespan is up to 40 years in captivity.
Wild asses can run swiftly, almost as fast as a horse. However, unlike most hoofed mammals, their tendency is to not flee right away from a potentially dangerous situation, but to investigate first before deciding what to do. When they need to, they can defend themselves with kicks from both their front and hind legs. Equids were used in ancient Sumer to pull wagons circa 2600 BC, and then chariots on the Standard of Ur, circa 2000 BC. These have been suggested to represent onagers, but are now thought to have been domesticated asses.
Diet
The African wild asses' diet consists of grasses, bark, and leaves. Despite being primarily adapted for living in an arid climate, they are dependent on water, and when not receiving the needed moisture from vegetation, they must drink at least once every three days. However, they can survive on a surprisingly small amount of liquid, and have been reported to drink salty or brackish water.
Conservation status
Though the species itself is under no threat of extinction, due to abundant domestic stock (donkeys and burros), the two extant wild subspecies are both listed as critically endangered. African wild asses have been captured for domestication for centuries and this, along with interbreeding between wild and domestic animals, has caused a distinct decline in population numbers. There are now only a few hundred individuals left in the wild. These animals are also hunted for food and for traditional medicine in both Ethiopia and Somalia. Competition with domestic livestock for grazing, and restricted access to water supplies caused by agricultural developments, pose further threats to the survival of this species. The African wild ass is legally protected in the countries where it is currently found, although these measures often prove difficult to enforce. A protected population of the Somali wild ass exists in the Yotvata Hai-Bar Nature Reserve in Israel, to the north of Eilat. This reserve was established in 1968 with the view to bolster populations of endangered desert species. Populations of horses and asses are fairly resilient and, if the species is properly protected, it may well recover from its current low.
In captivity
There are about 150 individual Somali wild asses living in zoos around the globe, of which 36 were born at Zoo Basel, where this species' breeding program started with Basel's first Somali wild asses in 1970 and the first birth in 1972.
Zoo Basel manages the European studbook for the Somali wild ass and coordinates the European Endangered Species Programme (EEP). All European wild donkeys are either descendants of the original group at Zoo Basel or of 12 others that came from the Yotvata Hai-Bar Nature Reserve in Israel in 1972.
| Biology and health sciences | Equidae | Animals |
2736310 | https://en.wikipedia.org/wiki/Domestic%20goose | Domestic goose | A domestic goose is a goose that humans have domesticated and kept for their meat, eggs, or down feathers, or as companion animals. Domestic geese have been derived through selective breeding from the wild greylag goose (Anser anser domesticus) and swan goose (Anser cygnoides domesticus).
Origins
In Europe, northern Africa, and western Asia, the original domesticated geese are derived from the greylag goose (Anser anser). In eastern Asia, the original domesticated geese are derived from the swan goose (Anser cygnoides); these are commonly known as Chinese geese. Both have been widely introduced in more recent times, and modern flocks in both areas (and elsewhere, such as Australia and North America) may consist of either species or hybrids between them. Chinese geese may be readily distinguished from European geese by the large knob at the base of the bill, though hybrids may exhibit every degree of variation between the two species.
Charles Darwin remarked in The Variation of Animals and Plants Under Domestication that the domestication of geese is of a very ancient date. There is archaeological evidence for domesticated geese in Egypt more than 4,000 years ago. It has been proposed that geese were domesticated around 3000 BCE in southeastern Europe, possibly in Greece, but reliable evidence of domestic geese comes from a much later period (8th century BCE) in the Odyssey. Another potential domestication site is in Egypt during the Old Kingdom (2686–1991 BCE) due to iconographic evidence of goose exploitation, but this scenario for the original domestication event has been considered less likely. Geese were also herded by ancient Mesopotamians for food and sacrifices and depicted in Mesopotamian art from the early Dynastic Period (2900–2350 BCE) onwards. Certainly, fully domesticated geese were present during the New Kingdom times in Egypt (1552–1151 BCE) and contemporaneously in Europe, and goose husbandry involving several varieties was well established by the Romans by the 1st century BCE. In the Medieval Period, goose husbandry was at its peak with large flocks kept by peasants. Archaeological evidence of the domestic goose in northern Europe indicates that it was probably introduced into Scandinavia during the Early Iron Age (400 BCE–550 CE).
Characteristics
Domestic geese have been selectively bred for size, with some breeds weighing up to , compared to the maximum of for the wild swan goose and for the wild greylag goose. This affects their body structure; whereas wild geese have a horizontal posture and slim rear end, domesticated geese lay down large fat deposits toward the tail end, giving a fat rear and forcing the bird into a more upright posture. Although their heavy weight affects their ability to fly, most breeds of domestic geese are capable of flight.
Geese have also been strongly selected for fecundity, with females laying up to 500 eggs per year, compared to 5–12 eggs for a wild goose.
As most domestic geese display little sexual dimorphism, sexing is based primarily on physical characteristics and behaviour. Males are typically taller and larger than females, and have longer, thicker necks. In addition, males can be distinguished by the protective behaviour they exhibit towards their mates and their offspring – the male will typically stand between his partner and any perceived threat.
Changes to the plumage are variable; many have been selected to lose dark brown tones of the wild bird. The result is an animal marked, or completely covered in white feathers. Others retain plumage close to the natural; some, such as the modern Toulouse goose look almost identical to the greylag in plumage, differing only in structure. White geese are often preferred as they look better plucked and dressed, with any small down feathers remaining being less conspicuous. From the time of the Romans, white geese have been held in great esteem.
Geese produce large edible eggs, weighing . They can be used in cooking just like chicken's eggs, though they have proportionally more yolk, and this cooks to a slightly denser consistency. The taste is much the same as that of a chicken egg, but gamier.
Like their wild ancestors, domestic geese are very protective of their offspring and other members of the flock. The gander will normally place himself between any perceived threat and his family. Owing to their highly aggressive nature, loud call and sensitivity to unusual movements, geese can contribute towards the security of a property. In late 1950s South Vietnam, the VNAF used flocks of geese to guard their parked aircraft at night due to the noise they would make at intruders.
Because domestic geese descended from the greylag goose are effectively the same species as their wild ancestor (being a subspecies formed through domestication), escaped individuals readily breed with wild populations, resulting in the offspring sometimes resembling either one of their parents, or bearing mixed plumage with patterns of grey and white feathers and orange beaks.
Due to their tendency to make noise when approached by strangers, about 500 geese were used to supplement dogs, drones, and humans to patrol the 533-km boundary between Chongzuo and Vietnam during the COVID-19 pandemic. An official commented that the birds, one of the most common livestock in the region, are sensitive to sounds and can sometimes be more aggressive than dogs.
Culinary uses
Geese are important to multiple culinary traditions. The meat, liver and other organs, fat, skin, blood and eggs are used culinarily in various cuisines.
Geese in fiction and myth
When Aphrodite first came ashore she was welcomed by the Charites (Roman "Graces"), whose chariot was drawn by geese.
There are Mother Goose tales, such as a farmwife might have told; there is the proverbial goose that laid the golden eggs, warning about the perils of greed.
The geese in the temple of Juno on the Capitoline Hill were said by Livy to have saved Rome from the Gauls around 390 BC when they were disturbed in a night attack. The story may be an attempt to explain the origin of the sacred flock of geese at Rome.
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| Biology and health sciences | Anseriformes | Animals |
2737484 | https://en.wikipedia.org/wiki/English%20landscape%20garden | English landscape garden | The English landscape garden, also called English landscape park or simply the English garden (, , , , ), is a style of "landscape" garden which emerged in England in the early 18th century, and spread across Europe, replacing the more formal, symmetrical French formal garden which had emerged in the 17th century as the principal gardening style of Europe. The English garden presented an idealized view of nature. Created and pioneered by William Kent and others, the "informal" garden style originated as a revolt against the architectural garden and drew inspiration from landscape paintings by Salvator Rosa, Claude Lorrain, and Nicolas Poussin, as well as from the classic Chinese gardens of the East, which had recently been described by European travellers and were realized in the Anglo-Chinese garden.
The English garden usually included a lake, sweeps of gently rolling lawns set against groves of trees, and recreations of classical temples, Gothic ruins, bridges, and other picturesque architecture, designed to recreate an idyllic pastoral landscape. The work of Lancelot "Capability" Brown was particularly influential. By the end of the 18th century the English garden was being imitated by the French landscape garden, and as far away as St. Petersburg, Russia, in Pavlovsk, the gardens of the future Emperor Paul. It also had a major influence on the forms of public parks and gardens which appeared around the world in the 19th century. The English landscape garden was usually centred on the English country house, and many examples in the United Kingdom are popular visitor attractions today.
History
The predecessors of the landscape garden in England were the great parks created by Sir John Vanbrugh (1664–1726) and Nicholas Hawksmoor at Castle Howard (1699–1712), Blenheim Palace (1705–1722), and the Claremont Landscape Garden at Claremont House (1715–1727). These parks featured vast lawns, woods, and pieces of architecture, such as the classical mausoleum designed by Hawksmoor at Castle Howard. At the centre of the composition was the house, behind which were formal and symmetrical gardens in the style of the garden à la française, with ornate carpets of floral designs and walls of hedges, decorated with statues and fountains. These gardens, modelled after the gardens of Versailles, were designed to impress visitors with their size and grandeur.
William Kent and Charles Bridgeman
The new style that became known as the English garden was invented by landscape designers William Kent and Charles Bridgeman, working for wealthy patrons, including Richard Temple, 1st Viscount Cobham; Richard Boyle, 3rd Earl of Burlington; and banker Henry Hoare. These men had large country estates, were members of the anti-royalist Whig Party, had classical educations, were patrons of the arts, and had taken the Grand Tour to Italy, where they had seen the Roman ruins and Italian landscapes they reproduced in their gardens.
William Kent (1685–1748) was an architect, painter and furniture designer who introduced Palladian-style architecture to England. Kent's inspiration came from Palladio's buildings in the Veneto and the landscapes and ruins around Rome – he lived in Italy from 1709 to 1719, and brought back many drawings of antique architecture and landscapes. His gardens were designed to complement the Palladian architecture of the houses he built.
Charles Bridgeman (1690–1738) was the son of a gardener and an experienced horticulturist, who became the Royal Gardener for Queen Anne and Prince George of Denmark, responsible for tending and redesigning the royal gardens at Windsor, Kensington Palace, Hampton Court, St. James's Park and Hyde Park. He collaborated with Kent on several major gardens, providing the botanical expertise which allowed Kent to realize his architectural visions.
Chiswick House
Kent created one of the first true English landscape gardens at Chiswick House for Richard Boyle, 3rd Earl of Burlington. The first gardens that he laid out between 1724 and 1733 had many formal elements of a garden à la française, including alleys forming a patte d'oie and canals, but they also featured a folly, a picturesque recreation of an Ionic temple set in a theatre of trees. Between 1733 and 1736, he redesigned the garden, adding lawns sloping down to the edge of the river and a small cascade. For the first time the form of a garden was inspired not by architecture, but by an idealized version of nature.
Rousham
Rousham House in Oxfordshire is considered by some as the most accomplished and significant of William Kent's work. The patron was General James Dormer, who commissioned Bridgeman to begin the garden in 1727, then brought in Kent to recreate it in 1737. Bridgeman had built a series of garden features including a grotto of Venus on the slope along the River Cherwell, connected by straight alleys. Kent turned the alleys into winding paths, built a gently turning stream, used the natural landscape features and slopes, and created a series of views and tableaux decorated with allegorical statues of Apollo, a wounded gladiator, a lion attacking a horse, and other subjects. He placed eyecatchers, pieces of classical architecture, to decorate the landscape, and made use of the ha-ha, a concealed ditch that kept grazing animals out of the garden while giving an uninterrupted vista from within. Finally, he added cascades modelled on those of the garden of Villa Aldobrandini and Villa di Pratolino in Italy, to add movement and drama.
Stowe House
Stowe Gardens, in Buckinghamshire, (1730–1738), was an even more radical departure from the formal French garden. In the early 18th century, Richard Temple, 1st Viscount Cobham, had commissioned Charles Bridgeman to design a formal garden, with architectural decorations by John Vanbrugh. Bridgeman's design included an octagonal lake and a rotunda (1720–21) designed by Vanbrugh.
In the 1730s, William Kent and James Gibbs were appointed to work with Bridgeman, who died in 1738. Kent remade the lake in a more natural shape, and created a new kind of garden, which took visitors on a tour of picturesque landscapes. It eventually included a Palladian bridge (1738); a Temple of Venus (1731) in the form of a Palladian villa; a Temple of Ancient Virtues (1737), with statues of famous Greeks and Romans; a Temple of British Worthies (1734–1735), with statues of British heroes; and a Temple of Modern Virtues, which was deliberately left in ruins, which contained a headless statue of Robert Walpole, Cobham's political rival.
The garden attracted visitors from all over Europe, including Jean-Jacques Rousseau. It became the inspiration for landscape gardens in Britain and on the Continent.
Stourhead
Stourhead, in Wiltshire (1741–1780), created by banker Henry Hoare, was one of the first 'picturesque' gardens, inspired to resemble the paintings of Claude Lorrain. Hoare had travelled to Italy on the Grand Tour and had returned with a painting by Claude Lorrain. Hoare dammed a stream on his estate, created a lake, and surrounded the lake with landscapes and architectural constructions representing the different steps of the journey of Aeneas in the Aeneid by Virgil.
The great age of the English garden
Capability Brown
The most influential figure in the later development of the English landscape garden was Lancelot "Capability" Brown (1716–1783), who began his career in 1740 as a gardener at Stowe Gardens under Charles Bridgeman, then succeeded William Kent in 1748.
Brown's contribution was to simplify the garden by eliminating geometric structures, alleys, and parterres near the house and replacing them with rolling lawns and extensive views out to isolated groups of trees, making the landscape seem even larger. "He sought to create an ideal landscape out of the English countryside." He created artificial lakes and used dams and canals to transform streams or springs into the illusion that a river flowed through the garden.
He compared his own role as a garden designer to that of a poet or composer. "Here I put a comma, there, when it's necessary to cut the view, I put a parenthesis; there I end it with a period and start on another theme."
Brown designed 170 gardens. The most important were:
Petworth (West Sussex) in 1752;
Chatsworth (Derbyshire) in 1761;
Bowood (Wiltshire) in 1763;
Blenheim Palace (Oxfordshire) in 1764.
Humphry Repton
Humphry Repton (21 April 1752 – 24 March 1818) was the last great English landscape designer of the eighteenth century, often regarded as the successor to Capability Brown. Repton hit upon the idea of becoming a "landscape gardener" (a term he himself coined) after failing at various ventures and, sensing an opportunity after Brown's death, was ambitious to fill the gap and sent circulars round his contacts in the upper classes advertising his services. To help clients visualize his designs, Repton produced 'Red Books' (so called for their binding) with explanatory text and watercolors with a system of overlays to show 'before' and 'after' views.
In 1794 Richard Payne Knight and Uvedale Price simultaneously published vicious attacks on the 'meagre genius of the bare and bald', criticizing Brown's smooth, serpentine curves as bland and unnatural and championing rugged and intricate designs, composed according to 'picturesque theory' that designed landscapes should be composed like landscape paintings, with a foreground, a middle ground and a background. Early in his career, Repton defended Brown's reputation during the 'picturesque controversy'. However, as his career progressed Repton came to apply picturesque theory to the practice of landscape design. He believed that the foreground should be the realm of art (with formal geometry and ornamental planting), that the middle ground should have a parkland character of the type created by Brown and that the background should have a wild and 'natural' character. Repton re-introduced formal terraces, balustrades, trellis work and flower gardens around the house in a way that became common practice in the nineteenth century.
Repton published four major books on garden design: Sketches and Hints on Landscape Gardening (1795), Observations on the Theory and Practice of Landscape Gardening (1803), An Inquiry into the Changes of Taste in Landscape Gardening (1806) and Fragments on the Theory and Practice of Landscape Gardening (1816). These drew on material and techniques used in the Red Books. These works greatly influenced other landscape-designers including John Claudius Loudon, John Nash, Jean-Charles Adolphe Alphand, Hermann Ludwig Heinrich Pückler-Muskau and Frederick Law Olmsted.
The "forest or savage garden"
One aspect of the new style was making woodland more interesting and ornamental, leading to the establishment of the woodland garden as a distinct type. This took several forms, one of which was helped by the developing Gothic revival. Horace Walpole, a great promoter of the English landscape garden style, praised Painshill in Surrey, whose varied features included a shrubbery with American plants, and a sloping "Alpine Valley" of conifers, as one of the best of the new style of "forest or savage gardens". This was a style of woodland aiming at the sublime, a newly-fashionable concept in literature and the arts, or at the least to be picturesque, another new term. It really required steep slopes, even if not very high, along which paths could be made revealing dramatic views, by which contemporary viewers who had read Gothic novels like Walpole's The Castle of Otranto (1764) were very ready to be impressed.
The appropriate style of garden buildings was Gothic rather than Neoclassical, and exotic planting was more likely to be evergreen conifers rather than flowering plants, replacing "the charm of bright, pleasant scenery in favour of the dark and rugged, gloomy and dramatic". A leading example of the style was Studley Royal in North Yorkshire, which had the great advantage, at what was known as "The Surprise View", of suddenly revealing a distant view from above of the impressive ruins of Fountains Abbey.
At Stowe, Capability Brown followed the new fashion between 1740 and 1753 by adding a new section to the park, called Hawkwelle Hill or the Gothic promenade, with a Gothic revival building. Walpole had decided in 1751 "to go Gothic", as he put it in a letter, and thereafter was a leading propagandist for the style, with his own house, Strawberry Hill in Twickenham, still the most extreme example of 18th-century "Gothick" style.
The "Anglo-Chinese" garden
According to some writers, especially French ones, the Far East inspired the origins of the English landscape garden, via Holland. In 1685, the English writer, formerly a diplomat at The Hague, Sir William Temple wrote an essay Upon the garden of Epicurus (published in 1690), including a passage which contrasted European symmetrical and formal gardens with asymmetrical compositions from China, for which he introduced (as Chinese) the term sharawadgi, in fact probably a mangled Japanese word for "irregularity". Temple had never visited the Far East, but he was in contact with the Dutch and their discourse on irregularity in design, had spoken to a merchant who had been in the Far East for a long time, and read the works of European travellers there. He noted that Chinese gardens avoided formal rows of trees and flower beds, and instead placed trees, plants, and other garden features in irregular ways to strike the eye and create beautiful compositions, with an understatement criticizing the formal compositions of the gardens at the Palace of Versailles of Louis XIV of France. His observations on the Chinese garden were cited by the essayist Joseph Addison in an essay in 1712, who used them to attack the English gardeners who, instead of imitating nature, tried to make their gardens in the French style, as far from nature as possible.
The novelty and exoticism of Chinese art and architecture in Europe led in 1738 to the construction of the first Chinese-style building in an English garden, in the garden of Stowe House, at a time when chinoiserie was popular in most forms of the decorative arts across Europe. The style became even more popular thanks to William Chambers (1723–1796), who lived in China from 1745 to 1747, and wrote a book, Designs of Chinese Buildings, Furniture, Dresses, Machines, and Utensils. To which is annexed, a Description of their Temples, Houses, Gardens, &c. published in 1757. In 1761 he built the Great Pagoda, London, as part of Kew Gardens, a park with gardens and architecture symbolizing all parts of the world and all architectural styles. Thereafter Chinese pagodas began to appear in other English gardens, then in France and elsewhere on the continent. French observers coined the term Jardin Anglo-Chinois (Anglo-Chinese garden) for this style of garden.
The English garden spreads to Continental Europe
Descriptions of English gardens were first brought to France by Jean-Bernard, abbé Le Blanc, who published accounts of his voyage in 1745 and 1751. A treatise, and tour guide, on the English garden, Observations on Modern Gardening, written by Thomas Whately and published in London in 1770, was translated into French and German in 1771. After the end of the Seven Years' War in 1763, French noblemen were able to voyage to England and see the gardens for themselves, and the style began to be adapted in French gardens. The new style also had the advantage of requiring fewer gardeners, and was easier to maintain, than the French garden.
One of the first English gardens on the continent was at Ermenonville, in France, built by marquis René Louis de Girardin from 1763 to 1776 and based on the ideals of Jean Jacques Rousseau, who was buried within the park. Rousseau and the garden's founder had visited Stowe a few years earlier. Other early examples were the Désert de Retz, Yvelines (1774–1782); the Gardens of the Château de Bagatelle in the Bois de Boulogne, west of Paris (1777–1784);
The Folie Saint James, in Neuilly-sur-Seine, (1777–1780); and the Château de Méréville, in the Essonne department, (1784–1786). Even at Versailles, the home of the most classical of all French gardens, a small English landscape park with a Roman temple was built and a mock village, the Hameau de la Reine (1783–1789), was created for Marie Antoinette.
The new style also spread to Germany. The central English Grounds of Wörlitz, in the Principality of Anhalt, was laid out between 1769 and 1773 by Leopold III, Duke of Anhalt-Dessau, based on the models of Claremont, Stourhead and Stowe Landscape Gardens. Another notable example was The Englischer Garten in Munich, Germany, created in 1789 by Sir Benjamin Thompson (1753–1814).
In the Netherlands the landscape-architect Lucas Pieters Roodbaard (1782–1851) designed several gardens and parks in this style. The style was introduced to Sweden by Fredrik Magnus Piper.
In Poland the main example of this style is Łazienki Park in Warsaw. The garden scheme owes its shape and appearance mainly to the last king of the country Stanisław August Poniatowski. In another part of the Polish–Lithuanian Commonwealth the Sofiyivka Park (Zofiówka), now Ukraine, was designed by Count Stanisław Szczęsny Potocki so as to illustrate the Odyssey and the Iliad.
The style also spread rapidly to Russia, where in 1774 Catherine the Great adapted the new style in the park of her palace at Tsarskoe Selo, complete with a mock Chinese village and a Palladian bridge, modeled after that at Wilton House. A much larger park was created for her son Paul in the neighbouring estate of Pavlovsk. The Monrepos Park is sited on the rocky island of Linnasaari in the Vyborg Bay and is noted for its glacially deposited boulders and granite rocks.
Characteristics of the English garden abroad
The continental European "English garden" is characteristically on a smaller scale; many are in or on the edge of cities, rather than in the middle of the countryside. Such gardens usually lack the sweeping vistas of gently rolling ground and water, which in England tend to be set against a woodland background with clumps of trees and outlier groves. Instead, they are often more densely studded with "eye-catchers", such as grottoes, temples, tea-houses, belvederes, pavilions, sham ruins, bridges, and statues. The name English garden – not used in the United Kingdom, where "landscape garden" serves – differentiates it from the formal Baroque design of the garden à la française. One of the best-known English gardens in Europe is the Englischer Garten in Munich.
The dominant style was revised in the early 19th century to include more "gardenesque" features, including shrubberies with gravelled walks, tree plantations to satisfy botanical curiosity, and, most notably, the return of flowers, in skirts of sweeping planted beds. This is the version of the landscape garden most imitated in Europe in the 19th century. The outer areas of the "home park" of English country houses retain their naturalistic shaping. English gardening since the 1840s has been on a more restricted scale, closer and more allied to the residence.
The canonical European English park contains a number of Romantic elements. Always present is a pond or small lake with a pier or bridge. Overlooking the pond is a round or hexagonal pavilion, often in the shape of a monopteros, a Roman temple. Sometimes the park also has a "Chinese" pavilion. Other elements include a grotto and imitation ruins.
A second style of English garden, which became popular during the 20th century in France and northern Europe, is based on the style of the late 19th-century English cottage garden, with abundant mixed planting of flowers, intended to appear largely unplanned.
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21171721 | https://en.wikipedia.org/wiki/Sea%20level%20rise | Sea level rise | Between 1901 and 2018, the average sea level rose by , with an increase of per year since the 1970s. This was faster than the sea level had ever risen over at least the past 3,000 years. The rate accelerated to /yr for the decade 2013–2022. Climate change due to human activities is the main cause. Between 1993 and 2018, melting ice sheets and glaciers accounted for 44% of sea level rise, with another 42% resulting from thermal expansion of water.
Sea level rise lags behind changes in the Earth's temperature by many decades, and sea level rise will therefore continue to accelerate between now and 2050 in response to warming that has already happened. What happens after that depends on human greenhouse gas emissions. If there are very deep cuts in emissions, sea level rise would slow between 2050 and 2100. It could then reach by 2100 slightly over from now and approximately from the 19th century. With high emissions it would instead accelerate further, and could rise by or even by 2100. In the long run, sea level rise would amount to over the next 2000 years if warming stays to its current over the pre-industrial past. It would be if warming peaks at .
Rising seas affect every coastal and island population on Earth. This can be through flooding, higher storm surges, king tides, and tsunamis. There are many knock-on effects. They lead to loss of coastal ecosystems like mangroves. Crop yields may reduce because of increasing salt levels in irrigation water. Damage to ports disrupts sea trade. The sea level rise projected by 2050 will expose places currently inhabited by tens of millions of people to annual flooding. Without a sharp reduction in greenhouse gas emissions, this may increase to hundreds of millions in the latter decades of the century.
Local factors like tidal range or land subsidence will greatly affect the severity of impacts. For instance, sea level rise in the United States is likely to be two to three times greater than the global average by the end of the century. Yet, of the 20 countries with the greatest exposure to sea level rise, twelve are in Asia, including Indonesia, Bangladesh and the Philippines. The resilience and adaptive capacity of ecosystems and countries also varies, which will result in more or less pronounced impacts. The greatest impact on human populations in the near term will occur in the low-lying Caribbean and Pacific islands. Sea level rise will make many of them uninhabitable later this century.
Societies can adapt to sea level rise in multiple ways. Managed retreat, accommodating coastal change, or protecting against sea level rise through hard-construction practices like seawalls are hard approaches. There are also soft approaches such as dune rehabilitation and beach nourishment. Sometimes these adaptation strategies go hand in hand. At other times choices must be made among different strategies. Poorer nations may also struggle to implement the same approaches to adapt to sea level rise as richer states.
Observations
Between 1901 and 2018, the global mean sea level rose by about . More precise data gathered from satellite radar measurements found an increase of from 1993 to 2017 (average of /yr). This accelerated to /yr for 2013–2022. Paleoclimate data shows that this rate of sea level rise is the fastest it had been over at least the past 3,000 years.
While sea level rise is uniform around the globe, some land masses are moving up or down as a consequence of subsidence (land sinking or settling) or post-glacial rebound (land rising as melting ice reduces weight). Therefore, local relative sea level rise may be higher or lower than the global average. Changing ice masses also affect the distribution of sea water around the globe through gravity.
Projections
Approaches used for projections
Several approaches are used for sea level rise (SLR) projections. One is process-based modeling, where ice melting is computed through an ice-sheet model and rising sea temperature and expansion through a general circulation model, and then these contributions are added up. The so-called semi-empirical approach instead applies statistical techniques and basic physical modeling to the observed sea level rise and its reconstructions from the historical geological data (known as paleoclimate modeling). It was developed because process-based model projections in the past IPCC reports (such as the Fourth Assessment Report from 2007) were found to underestimate the already observed sea level rise.
By 2013, improvements in modeling had addressed this issue, and model and semi-empirical projections for the year 2100 are now very similar. Yet, semi-empirical estimates are reliant on the quality of available observations and struggle to represent non-linearities, while processes without enough available information about them cannot be modeled. Thus, another approach is to combine the opinions of a large number of scientists in what is known as a structured expert judgement (SEJ).
Variations of these primary approaches exist. For instance, large climate models are always in demand, so less complex models are often used in their place for simpler tasks like projecting flood risk in the specific regions. A structured expert judgement may be used in combination with modeling to determine which outcomes are more or less likely, which is known as "shifted SEJ". Semi-empirical techniques can be combined with the so-called "intermediate-complexity" models. After 2016, some ice sheet modeling exhibited the so-called ice cliff instability in Antarctica, which results in substantially faster disintegration and retreat than otherwise simulated. The differences are limited with low warming, but at higher warming levels, ice cliff instability predicts far greater sea level rise than any other approach.
Projections for the 21st century
The Intergovernmental Panel on Climate Change is the largest and most influential scientific organization on climate change, and since 1990, it provides several plausible scenarios of 21st century sea level rise in each of its major reports. The differences between scenarios are mainly due to uncertainty about future greenhouse gas emissions. These depend on future economic developments, and also future political action which is hard to predict. Each scenario provides an estimate for sea level rise as a range with a lower and upper limit to reflect the unknowns. The scenarios in the 2013–2014 Fifth Assessment Report (AR5) were called Representative Concentration Pathways, or RCPs and the scenarios in the IPCC Sixth Assessment Report (AR6) are known as Shared Socioeconomic Pathways, or SSPs. A large difference between the two was the addition of SSP1-1.9 to AR6, which represents meeting the best Paris climate agreement goal of . In that case, the likely range of sea level rise by 2100 is .
The lowest scenario in AR5, RCP2.6, would see greenhouse gas emissions low enough to meet the goal of limiting warming by 2100 to . It shows sea level rise in 2100 of about with a range of . The "moderate" scenario, where emissions take a decade or two to peak and its atmospheric concentration does not plateau until the 2070s is called RCP 4.5. Its likely range of sea level rise is . The highest scenario in RCP8.5 pathway sea level would rise between . AR6 had equivalents for both scenarios, but it estimated larger sea level rise under both. In AR6, the SSP1-2.6 pathway results in a range of by 2100. The "moderate" SSP2-4.5 results in a range by 2100 and SSP5-8.5 led to .
This general increase of projections in AR6 came after the improvements in ice-sheet modeling and the incorporation of structured expert judgements. These decisions came as the observed ice-sheet erosion in Greenland and Antarctica had matched the upper-end range of the AR5 projections by 2020, and the finding that AR5 projections were likely too slow next to an extrapolation of observed sea level rise trends, while the subsequent reports had improved in this regard. Further, AR5 was criticized by multiple researchers for excluding detailed estimates the impact of "low-confidence" processes like marine ice sheet and marine ice cliff instability, which can substantially accelerate ice loss to potentially add "tens of centimeters" to sea level rise within this century. AR6 includes a version of SSP5-8.5 where these processes take place, and in that case, sea level rise of up to by 2100 could not be ruled out.
Role of instability processes
The greatest uncertainty with sea level rise projections is associated with the so-called marine ice sheet instability (MISI), and, even more so, Marine Ice Cliff Instability (MICI). These processes are mainly associated with West Antarctic Ice Sheet, but may also apply to some of Greenland's glaciers. The former suggests that when glaciers are mostly underwater on retrograde (backwards-sloping) bedrock, the water melts more and more of their height as their retreat continues, thus accelerating their breakdown on its own. This is widely accepted, but is difficult to model.
The latter posits that coastal ice cliffs which exceed ~ in above-ground height and are ~ in basal (underground) height are likely to rapidly collapse under their own weight once the ice shelves propping them up are gone. The collapse then exposes the ice masses following them to the same instability, potentially resulting in a self-sustaining cycle of cliff collapse and rapid ice sheet retreat. This theory had been highly influential - in a 2020 survey of 106 experts, the 2016 paper which suggested or more of sea level rise by 2100 from Antarctica alone, was considered even more important than the 2014 IPCC Fifth Assessment Report. Even more rapid sea level rise was proposed in a 2016 study led by Jim Hansen, which hypothesized multi-meter sea level rise in 50–100 years as a plausible outcome of high emissions, but it remains a minority view amongst the scientific community.
Marine ice cliff instability had also been very controversial, since it was proposed as a modelling exercise, and the observational evidence from both the past and the present is very limited and ambiguous. So far, only one episode of seabed gouging by ice from the Younger Dryas period appears truly consistent with this theory, but it had lasted for an estimated 900 years, so it is unclear if it supports rapid sea level rise in the present. Modelling which investigated the hypothesis after 2016 often suggested that the ice shelves in the real world may collapse too slowly to make this scenario relevant, or that ice mélange - debris produced as the glacier breaks down - would quickly build up in front of the glacier and significantly slow or even outright stop the instability soon after it began.
Due to these uncertainties, some scientists - including the originators of the hypothesis, Robert DeConto and David Pollard - have suggested that the best way to resolve the question would be to precisely determine sea level rise during the Last Interglacial. MICI can be effectively ruled out if SLR at the time was lower than , while it is very likely if the SLR was greater than . As of 2023, the most recent analysis indicates that the Last Interglacial SLR is unlikely to have been higher than , as higher values in other research, such as , appear inconsistent with the new paleoclimate data from The Bahamas and the known history of the Greenland Ice Sheet.
Post-2100 sea level rise
Even if the temperature stabilizes, significant sea-level rise (SLR) will continue for centuries, consistent with paleo records of sea level rise. This is due to the high level of inertia in the carbon cycle and the climate system, owing to factors such as the slow diffusion of heat into the deep ocean, leading to a longer climate response time. A 2018 paper estimated that sea level rise in 2300 would increase by a median of 20 cm (8 in) for every five years emissions increase before peaking. It shows a 5% likelihood of a increase due to the same. The same estimate found that if the temperature stabilized below , 2300 sea level rise would still exceed . Early net zero and slowly falling temperatures could limit it to .
By 2021, the IPCC Sixth Assessment Report was able to provide estimates for sea level rise in 2150. Keeping warming to 1.5°C under the SSP1-1.9 scenario would result in sea level rise in the 17–83% range of . In the SSP1-2.6 pathway the range would be , for SSP2-4.5 a range by 2100 and for SSP5-8.5 a rise of . It stated that the "low-confidence, high impact" projected mean sea level rise by 2100, and that by 2150, the total sea level rise in his scenario would be in the range of by 2150. AR6 also provided lower-confidence estimates for year 2300 sea level rise under SSP1-2.6 and SSP5-8.5 with various impact assumptions. In the best case scenario, under SSP1-2.6 with no ice sheet acceleration after 2100, the estimate was only . In the worst estimated scenario, SSP-8.5 with ice cliff instability, the projected range for total sea level rise was by the year 2300.
Projections for subsequent years are more difficult. In 2019, when 22 experts on ice sheets were asked to estimate 2200 and 2300 SLR under the 5°C warming scenario, there were 90% confidence intervals of − to and − to , respectively. (Negative values represent the extremely low probability of large climate change-induced increases in precipitation greatly elevating ice sheet surface mass balance.) In 2020, 106 experts who contributed to 6 or more papers on sea level estimated median SLR in the year 2300 for the low-warming RCP2.6 scenario and the median of for the high-warming RCP8.5. The former scenario had the 5%–95% confidence range of , and the latter of .
After 500 years, sea level rise from thermal expansion alone may have reached only half of its eventual level - likely within ranges of . Additionally, tipping points of Greenland and Antarctica ice sheets are likely to play a larger role over such timescales. Ice loss from Antarctica is likely to dominate very long-term SLR, especially if the warming exceeds . Continued carbon dioxide emissions from fossil fuel sources could cause additional tens of metres of sea level rise, over the next millennia. Burning of all fossil fuels on Earth is sufficient to melt the entire Antarctic ice sheet, causing about of sea level rise.
Year 2021 IPCC estimates for the amount of sea level rise over the next 2,000 years project that:
At a warming peak of , global sea levels would rise
At a warming peak of , sea levels would rise
At a warming peak of , sea levels would rise
Sea levels would continue to rise for several thousand years after the ceasing of emissions, due to the slow nature of climate response to heat. The same estimates on a timescale of 10,000 years project that:
At a warming peak of , global sea levels would rise
At a warming peak of , sea levels would rise
At a warming peak of , sea levels would rise
Measurements
Variations in the amount of water in the oceans, changes in its volume, or varying land elevation compared to the sea surface can drive sea level changes. Over a consistent time period, assessments can attribute contributions to sea level rise and provide early indications of change in trajectory. This helps to inform adaptation plans. The different techniques used to measure changes in sea level do not measure exactly the same level. Tide gauges can only measure relative sea level. Satellites can also measure absolute sea level changes. To get precise measurements for sea level, researchers studying the ice and oceans factor in ongoing deformations of the solid Earth. They look in particular at landmasses still rising from past ice masses retreating, and the Earth's gravity and rotation.
Satellites
Since the launch of TOPEX/Poseidon in 1992, an overlapping series of altimetric satellites has been continuously recording the sea level and its changes. These satellites can measure the hills and valleys in the sea caused by currents and detect trends in their height. To measure the distance to the sea surface, the satellites send a microwave pulse towards Earth and record the time it takes to return after reflecting off the ocean's surface. Microwave radiometers correct the additional delay caused by water vapor in the atmosphere. Combining these data with the location of the spacecraft determines the sea-surface height to within a few centimetres. These satellite measurements have estimated rates of sea level rise for 1993–2017 at per year.
Satellites are useful for measuring regional variations in sea level. An example is the substantial rise between 1993 and 2012 in the western tropical Pacific. This sharp rise has been linked to increasing trade winds. These occur when the Pacific Decadal Oscillation (PDO) and the El Niño–Southern Oscillation (ENSO) change from one state to the other. The PDO is a basin-wide climate pattern consisting of two phases, each commonly lasting 10 to 30 years. The ENSO has a shorter period of 2 to 7 years.
Tide gauges
The global network of tide gauges is the other important source of sea-level observations. Compared to the satellite record, this record has major spatial gaps but covers a much longer period. Coverage of tide gauges started mainly in the Northern Hemisphere. Data for the Southern Hemisphere remained scarce up to the 1970s. The longest running sea-level measurements, NAP or Amsterdam Ordnance Datum were established in 1675, in Amsterdam. Record collection is also extensive in Australia. They include measurements by Thomas Lempriere, an amateur meteorologist, beginning in 1837. Lempriere established a sea-level benchmark on a small cliff on the Isle of the Dead near the Port Arthur convict settlement in 1841.
Together with satellite data for the period after 1992, this network established that global mean sea level rose between 1870 and 2004 at an average rate of about 1.44 mm/yr. (For the 20th century the average is 1.7 mm/yr.) By 2018, data collected by Australia's Commonwealth Scientific and Industrial Research Organisation (CSIRO) had shown that the global mean sea level was rising by per year. This was double the average 20th century rate. The 2023 World Meteorological Organization report found further acceleration to 4.62 mm/yr over the 2013–2022 period. These observations help to check and verify predictions from climate change simulations.
Regional differences are also visible in the tide gauge data. Some are caused by local sea level differences. Others are due to vertical land movements. In Europe, only some land areas are rising while the others are sinking. Since 1970, most tidal stations have measured higher seas. However sea levels along the northern Baltic Sea have dropped due to post-glacial rebound.
Past sea level rise
An understanding of past sea level is an important guide to where current changes in sea level will end up. In the recent geological past, thermal expansion from increased temperatures and changes in land ice are the dominant reasons of sea level rise. The last time that the Earth was warmer than pre-industrial temperatures was 120,000 years ago. This was when warming due to Milankovitch cycles (changes in the amount of sunlight due to slow changes in the Earth's orbit) caused the Eemian interglacial. Sea levels during that warmer interglacial were at least higher than now. The Eemian warming was sustained over a period of thousands of years. The size of the rise in sea level implies a large contribution from the Antarctic and Greenland ice sheets. Levels of atmospheric carbon dioxide of around 400 parts per million (similar to 2000s) had increased temperature by over around three million years ago. This temperature increase eventually melted one third of Antarctica's ice sheet, causing sea levels to rise 20 meters above the preindustrial levels.
Since the Last Glacial Maximum, about 20,000 years ago, sea level has risen by more than . Rates vary from less than 1 mm/year during the pre-industrial era to 40+ mm/year when major ice sheets over Canada and Eurasia melted. Meltwater pulses are periods of fast sea level rise caused by the rapid disintegration of these ice sheets. The rate of sea level rise started to slow down about 8,200 years before today. Sea level was almost constant for the last 2,500 years. The recent trend of rising sea level started at the end of the 19th or beginning of the 20th century.
Causes
Effects of climate change
The three main reasons why global warming causes sea levels to rise are the expansion of oceans due to heating, water inflow from melting ice sheets and water inflow from glaciers. Other factors affecting sea level rise include changes in snow mass, and flow from terrestrial water storage, though the contribution from these is thought to be small. Glacier retreat and ocean expansion have dominated sea level rise since the start of the 20th century. Some of the losses from glaciers are offset when precipitation falls as snow, accumulates and over time forms glacial ice. If precipitation, surface processes and ice loss at the edge balance each other, sea level remains the same. Because of this precipitation began as water vapor evaporated from the ocean surface, effects of climate change on the water cycle can even increase ice build-up. However, this effect is not enough to fully offset ice losses, and sea level rise continues to accelerate.
The contributions of the two large ice sheets, in Greenland and Antarctica, are likely to increase in the 21st century. They store most of the land ice (~99.5%) and have a sea-level equivalent (SLE) of for Greenland and for Antarctica. Thus, melting of all the ice on Earth would result in about of sea level rise, although this would require at least 10,000 years and up to of global warming.
Ocean heating
The oceans store more than 90% of the extra heat added to the climate system by Earth's energy imbalance and act as a buffer against its effects. This means that the same amount of heat that would increase the average world ocean temperature by would increase atmospheric temperature by approximately . So a small change in the mean temperature of the ocean represents a very large change in the total heat content of the climate system. Winds and currents move heat into deeper parts of the ocean. Some of it reaches depths of more than .
When the ocean gains heat, the water expands and sea level rises. Warmer water and water under great pressure (due to depth) expand more than cooler water and water under less pressure. Consequently, cold Arctic Ocean water will expand less than warm tropical water. Different climate models present slightly different patterns of ocean heating. So their projections do not agree fully on how much ocean heating contributes to sea level rise.
Ice loss on the Antarctic continent
The large volume of ice on the Antarctic continent stores around 60% of the world's fresh water. Excluding groundwater this is 90%. Antarctica is experiencing ice loss from coastal glaciers in the West Antarctica and some glaciers of East Antarctica. However it is gaining mass from the increased snow build-up inland, particularly in the East. This leads to contradicting trends. There are different satellite methods for measuring ice mass and change. Combining them helps to reconcile the differences. However, there can still be variations between the studies. In 2018, a systematic review estimated average annual ice loss of 43 billion tons (Gt) across the entire continent between 1992 and 2002. This tripled to an annual average of 220 Gt from 2012 to 2017. However, a 2021 analysis of data from four different research satellite systems (Envisat, European Remote-Sensing Satellite, GRACE and GRACE-FO and ICESat) indicated annual mass loss of only about 12 Gt from 2012 to 2016. This was due to greater ice gain in East Antarctica than estimated earlier.
In the future, it is known that West Antarctica at least will continue to lose mass, and the likely future losses of sea ice and ice shelves, which block warmer currents from direct contact with the ice sheet, can accelerate declines even in East Antarctica. Altogether, Antarctica is the source of the largest uncertainty for future sea level projections. In 2019, the SROCC assessed several studies attempting to estimate 2300 sea level rise caused by ice loss in Antarctica alone, arriving at projected estimates of for the low emission RCP2.6 scenario, and in the high emission RCP8.5 scenario. This wide range of estimates is mainly due to the uncertainties regarding marine ice sheet and marine ice cliff instabilities.
East Antarctica
The world's largest potential source of sea level rise is the East Antarctic Ice Sheet (EAIS). It is 2.2 km thick on average and holds enough ice to raise global sea levels by 53.3 m (174 ft 10 in) Its great thickness and high elevation make it more stable than the other ice sheets. As of the early 2020s, most studies show that it is still gaining mass. Some analyses have suggested it began to lose mass in the 2000s. However they over-extrapolated some observed losses on to the poorly observed areas. A more complete observational record shows continued mass gain.
In spite of the net mass gain, some East Antarctica glaciers have lost ice in recent decades due to ocean warming and declining structural support from the local sea ice, such as Denman Glacier, and Totten Glacier. Totten Glacier is particularly important because it stabilizes the Aurora Subglacial Basin. Subglacial basins like Aurora and Wilkes Basin are major ice reservoirs together holding as much ice as all of West Antarctica. They are more vulnerable than the rest of East Antarctica. Their collective tipping point probably lies at around of global warming. It may be as high as or as low as . Once this tipping point is crossed, the collapse of these subglacial basins could take place over as little as 500 or as much as 10,000 years. The median timeline is 2000 years. Depending on how many subglacial basins are vulnerable, this causes sea level rise of between and .
On the other hand, the whole EAIS would not definitely collapse until global warming reaches , with a range between and . It would take at least 10,000 years to disappear. Some scientists have estimated that warming would have to reach at least to melt two thirds of its volume.
West Antarctica
East Antarctica contains the largest potential source of sea level rise. However the West Antarctic ice sheet (WAIS) is substantially more vulnerable. Temperatures on West Antarctica have increased significantly, unlike East Antarctica and the Antarctic Peninsula. The trend is between and per decade between 1976 and 2012. Satellite observations recorded a substantial increase in WAIS melting from 1992 to 2017. This resulted in of Antarctica sea level rise. Outflow glaciers in the Amundsen Sea Embayment played a disproportionate role.
The median estimated increase in sea level rise from Antarctica by 2100 is ~. There is no difference between scenarios, because the increased warming would intensify the water cycle and increase snowfall accumulation over the EAIS at about the same rate as it would increase ice loss from WAIS. However, most of the bedrock underlying the WAIS lies well below sea level, and it has to be buttressed by the Thwaites and Pine Island glaciers. If these glaciers were to collapse, the entire ice sheet would as well. Their disappearance would take at least several centuries, but is considered almost inevitable, as their bedrock topography deepens inland and becomes more vulnerable to meltwater, in what is known as marine ice sheet instability.
The contribution of these glaciers to global sea levels has already accelerated since the year 2000. The Thwaites Glacier now accounts for 4% of global sea level rise. It could start to lose even more ice if the Thwaites Ice Shelf fails and would no longer stabilize it, which could potentially occur in mid-2020s. A combination of ice sheet instability with other important but hard-to-model processes like hydrofracturing (meltwater collects atop the ice sheet, pools into fractures and forces them open) or smaller-scale changes in ocean circulation could cause the WAIS to contribute up to by 2100 under the low-emission scenario and up to under the highest-emission one. Ice cliff instability would cause a contribution of or more if it were applicable.
The melting of all the ice in West Antarctica would increase the total sea level rise to . However, mountain ice caps not in contact with water are less vulnerable than the majority of the ice sheet, which is located below the sea level. Its collapse would cause ~ of sea level rise. This disappearance would take an estimated 2000 years. The absolute minimum for the loss of West Antarctica ice is 500 years, and the potential maximum is 13,000 years.
Once ice loss from the West Antarctica is triggered, the only way to restore it to near-present values is by lowering the global temperature to below the preindustrial level. This would be below the temperature of 2020. Other researchers suggested that a climate engineering intervention to stabilize the ice sheet's glaciers may delay its loss by centuries and give more time to adapt. However this is an uncertain proposal, and would end up as one of the most expensive projects ever attempted.
Ice sheet loss in Greenland
Most ice on Greenland is in the Greenland ice sheet which is at its thickest. The rest of Greenland ice forms isolated glaciers and ice caps. The average annual ice loss in Greenland more than doubled in the early 21st century compared to the 20th century. Its contribution to sea level rise correspondingly increased from 0.07 mm per year between 1992 and 1997 to 0.68 mm per year between 2012 and 2017. Total ice loss from the Greenland ice sheet between 1992 and 2018 amounted to 3,902 gigatons (Gt) of ice. This is equivalent to a SLR contribution of 10.8 mm. The contribution for the 2012–2016 period was equivalent to 37% of sea level rise from land ice sources (excluding thermal expansion). This observed rate of ice sheet melting is at the higher end of predictions from past IPCC assessment reports.
In 2021, AR6 estimated that by 2100, the melting of Greenland ice sheet would most likely add around to sea levels under the low-emission scenario, and under the high-emission scenario. The first scenario, SSP1-2.6, largely fulfils the Paris Agreement goals, while the other, SSP5-8.5, has the emissions accelerate throughout the century. The uncertainty about ice sheet dynamics can affect both pathways. In the best-case scenario, ice sheet under SSP1-2.6 gains enough mass by 2100 through surface mass balance feedbacks to reduce the sea levels by . In the worst case, it adds . For SSP5-8.5, the best-case scenario is adding to sea levels, and the worst-case is adding .
Greenland's peripheral glaciers and ice caps crossed an irreversible tipping point around 1997. Sea level rise from their loss is now unstoppable. However the temperature changes in future, the warming of 2000–2019 had already damaged the ice sheet enough for it to eventually lose ~3.3% of its volume. This is leading to of future sea level rise. At a certain level of global warming, the Greenland ice sheet will almost completely melt. Ice cores show this happened at least once over the last million years, during which the temperatures have at most been warmer than the preindustrial average.
2012 modelling suggested that the tipping point of the ice sheet was between and . 2023 modelling has narrowed the tipping threshold to a - range, which is consistent with the empirical upper limit from ice cores. If temperatures reach or exceed that level, reducing the global temperature to above pre-industrial levels or lower would prevent the loss of the entire ice sheet. One way to do this in theory would be large-scale carbon dioxide removal, but there would still be cause of greater ice losses and sea level rise from Greenland than if the threshold was not breached in the first place. If the tipping point instead is durably but mildly crossed, the ice sheet would take between 10,000 and 15,000 years to disintegrate entirel, with a most likely estimate of 10,000 years. If climate change continues along its worst trajectory and temperatures continue to rise quickly over multiple centuries, it would only take 1,000 years.
Mountain glacier loss
There are roughly 200,000 glaciers on Earth, which are spread out across all continents. Less than 1% of glacier ice is in mountain glaciers, compared to 99% in Greenland and Antarctica. However, this small size also makes mountain glaciers more vulnerable to melting than the larger ice sheets. This means they have had a disproportionate contribution to historical sea level rise and are set to contribute a smaller, but still significant fraction of sea level rise in the 21st century. Observational and modelling studies of mass loss from glaciers and ice caps show they contribute 0.2-0.4 mm per year to sea level rise, averaged over the 20th century. The contribution for the 2012–2016 period was nearly as large as that of Greenland. It was 0.63 mm of sea level rise per year, equivalent to 34% of sea level rise from land ice sources. Glaciers contributed around 40% to sea level rise during the 20th century, with estimates for the 21st century of around 30%.
In 2023, a Science paper estimated that at , one quarter of mountain glacier mass would be lost by 2100 and nearly half would be lost at , contributing ~ and ~ to sea level rise, respectively. Glacier mass is disproportionately concentrated in the most resilient glaciers. So in practice this would remove 49-83% of glacier formations. It further estimated that the current likely trajectory of would result in the SLR contribution of ~ by 2100. Mountain glaciers are even more vulnerable over the longer term. In 2022, another Science paper estimated that almost no mountain glaciers could survive once warming crosses . Their complete loss is largely inevitable around . There is even a possibility of complete loss after 2100 at just . This could happen as early as 50 years after the tipping point is crossed, although 200 years is the most likely value, and the maximum is around 1000 years.
Sea ice loss
Sea ice loss contributes very slightly to global sea level rise. If the melt water from ice floating in the sea was exactly the same as sea water then, according to Archimedes' principle, no rise would occur. However melted sea ice contains less dissolved salt than sea water and is therefore less dense, with a slightly greater volume per unit of mass. If all floating ice shelves and icebergs were to melt sea level would only rise by about .
Changes to land water storage
Human activity impacts how much water is stored on land. Dams retain large quantities of water, which is stored on land rather than flowing into the sea, though the total quantity stored will vary from time to time. On the other hand, humans extract water from lakes, wetlands and underground reservoirs for drinking and food production. This often causes subsidence. Furthermore, the hydrological cycle is influenced by climate change and deforestation. In the 20th century, these processes had approximately cancelled out each other's impact on sea level rise, but dam building has slowed down and is expected to stay low for the 21st century.
Water redistribution caused by irrigation from 1993 to 2010 caused a drift of Earth's rotational pole by . This caused groundwater depletion equivalent to a global sea level rise of .
Impacts
On people and societies
Sea-level rise has many impacts. They include higher and more frequent high-tide and storm-surge flooding and increased coastal erosion. Other impacts are inhibition of primary production processes, more extensive coastal inundation, and changes in surface water quality and groundwater. These can lead to a greater loss of property and coastal habitats, loss of life during floods and loss of cultural resources. There are also impacts on agriculture and aquaculture. There can also be loss of tourism, recreation, and transport-related functions. Land use changes such as urbanisation or deforestation of low-lying coastal zones exacerbate coastal flooding impacts. Regions already vulnerable to rising sea level also struggle with coastal flooding. This washes away land and alters the landscape.
Changes in emissions are likely to have only a small effect on the extent of sea level rise by 2050. So projected sea level rise could put tens of millions of people at risk by then. Scientists estimate that 2050 levels of sea level rise would result in about 150 million people under the water line during high tide. About 300 million would be in places flooded every year. This projection is based on the distribution of population in 2010. It does not take into account the effects of population growth and human migration. These figures are 40 million and 50 million more respectively than the numbers at risk in 2010. By 2100, there would be another 40 million people under the water line during high tide if sea level rise remains low. This figure would be 80 million for a high estimate of median sea level rise. Ice sheet processes under the highest emission scenario would result in sea level rise of well over by 2100. This could be as much as over , This could result in as many as 520 million additional people ending up under the water line during high tide and 640 million in places flooded every year, compared to the 2010 population distribution.
Over the longer term, coastal areas are particularly vulnerable to rising sea levels. They are also vulnerable to changes in the frequency and intensity of storms, increased precipitation, and rising ocean temperatures. Ten percent of the world's population live in coastal areas that are less than above sea level. Two thirds of the world's cities with over five million people are located in these low-lying coastal areas. About 600 million people live directly on the coast around the world. Cities such as Miami, Rio de Janeiro, Osaka and Shanghai will be especially vulnerable later in the century under warming of 3 °C (5.4 °F). This is close to the current trajectory. LiDAR-based research had established in 2021 that 267 million people worldwide lived on land less than above sea level. With a sea level rise and zero population growth, that could increase to 410 million people.
Potential disruption of sea trade and migrations could impact people living further inland. United Nations Secretary-General António Guterres warned in 2023 that sea level rise risks causing human migrations on a "biblical scale". Sea level rise will inevitably affect ports, but there is limited research on this. There is insufficient knowledge about the investments necessary to protect ports currently in use. This includes protecting current facilities before it becomes more reasonable to build new ports elsewhere. Some coastal regions are rich agricultural lands. Their loss to the sea could cause food shortages. This is a particularly acute issue for river deltas such as Nile Delta in Egypt and Red River and Mekong Deltas in Vietnam. Saltwater intrusion into the soil and irrigation water has a disproportionate effect on them.
On ecosystems
Flooding and soil/water salinization threaten the habitats of coastal plants, birds, and freshwater/estuarine fish when seawater reaches inland. When coastal forest areas become inundated with saltwater to the point no trees can survive the resulting habitats are called ghost forests. Starting around 2050, some nesting sites in Florida, Cuba, Ecuador and the island of Sint Eustatius for leatherback, loggerhead, hawksbill, green and olive ridley turtles are expected to be flooded. The proportion will increase over time. In 2016, Bramble Cay islet in the Great Barrier Reef was inundated. This flooded the habitat of a rodent named Bramble Cay melomys. It was officially declared extinct in 2019.
Some ecosystems can move inland with the high-water mark. But natural or artificial barriers prevent many from migrating. This coastal narrowing is sometimes called 'coastal squeeze' when it involves human-made barriers. It could result in the loss of habitats such as mudflats and tidal marshes. Mangrove ecosystems on the mudflats of tropical coasts nurture high biodiversity. They are particularly vulnerable due to mangrove plants' reliance on breathing roots or pneumatophores. These will be submerged if the rate is too rapid for them to migrate upward. This would result in the loss of an ecosystem. Both mangroves and tidal marshes protect against storm surges, waves and tsunamis, so their loss makes the effects of sea level rise worse. Human activities such as dam building may restrict sediment supplies to wetlands. This would prevent natural adaptation processes. The loss of some tidal marshes is unavoidable as a consequence.
Corals are important for bird and fish life. They need to grow vertically to remain close to the sea surface in order to get enough energy from sunlight. The corals have so far been able to keep up the vertical growth with the rising seas, but might not be able to do so in the future.
Regional variations
When a glacier or ice sheet melts, it loses mass. This reduces its gravitational pull. In some places near current and former glaciers and ice sheets, this has caused water levels to drop. At the same time water levels will increase more than average further away from the ice sheet. Thus ice loss in Greenland affects regional sea level differently than the equivalent loss in Antarctica. On the other hand, the Atlantic is warming at a faster pace than the Pacific. This has consequences for Europe and the U.S. East Coast. The East Coast sea level is rising at 3–4 times the global average. Scientists have linked extreme regional sea level rise on the US Northeast Coast to the downturn of the Atlantic meridional overturning circulation (AMOC).
Many ports, urban conglomerations, and agricultural regions stand on river deltas. Here land subsidence contributes to much higher relative sea level rise. Unsustainable extraction of groundwater and oil and gas is one cause. Levees and other flood management practices are another. They prevent sediments from accumulating. These would otherwise compensate for the natural settling of deltaic soils.
Estimates for total human-caused subsidence in the Rhine-Meuse-Scheldt delta (Netherlands) are , over in urban areas of the Mississippi River Delta (New Orleans), and over in the Sacramento–San Joaquin River Delta. On the other hand, relative sea level around the Hudson Bay in Canada and the northern Baltic Sea is falling due to post-glacial isostatic rebound.
Adaptation
Cutting greenhouse gas emissions can slow and stabilize the rate of sea level rise after 2050. This would greatly reduce its costs and damages, but cannot stop it outright. So climate change adaptation to sea level rise is inevitable. The simplest approach is to stop development in vulnerable areas and ultimately move people and infrastructure away from them. Such retreat from sea level rise often results in the loss of livelihoods. The displacement of newly impoverished people could burden their new homes and accelerate social tensions.
It is possible to avoid or at least delay the retreat from sea level rise with enhanced protections. These include dams, levees or improved natural defenses. Other options include updating building standards to reduce damage from floods, addition of storm water valves to address more frequent and severe flooding at high tide, or cultivating crops more tolerant of saltwater in the soil, even at an increased cost. These options divide into hard and soft adaptation. Hard adaptation generally involves large-scale changes to human societies and ecological systems. It often includes the construction of capital-intensive infrastructure. Soft adaptation involves strengthening natural defenses and local community adaptation. This usually involves simple, modular and locally owned technology. The two types of adaptation may be complementary or mutually exclusive. Adaptation options often require significant investment. But the costs of doing nothing are far greater. One example would involve adaptation against flooding. Effective adaptation measures could reduce future annual costs of flooding in 136 of the world's largest coastal cities from $1 trillion by 2050 without adaptation to a little over $60 billion annually. The cost would be $50 billion per year. Some experts argue that retreat from the coast would have a lower impact on the GDP of India and Southeast Asia then attempting to protect every coastline, in the case of very high sea level rise.
To be successful, adaptation must anticipate sea level rise well ahead of time. As of 2023, the global state of adaptation planning is mixed. A survey of 253 planners from 49 countries found that 98% are aware of sea level rise projections, but 26% have not yet formally integrated them into their policy documents. Only around a third of respondents from Asian and South American countries have done so. This compares with 50% in Africa, and over 75% in Europe, Australasia and North America. Some 56% of all surveyed planners have plans which account for 2050 and 2100 sea level rise. But 53% use only a single projection rather than a range of two or three projections. Just 14% use four projections, including the one for "extreme" or "high-end" sea level rise. Another study found that over 75% of regional sea level rise assessments from the West and Northeastern United States included at least three estimates. These are usually RCP2.6, RCP4.5 and RCP8.5, and sometimes include extreme scenarios. But 88% of projections from the American South had only a single estimate. Similarly, no assessment from the South went beyond 2100. By contrast 14 assessments from the West went up to 2150, and three from the Northeast went to 2200. 56% of all localities were also found to underestimate the upper end of sea level rise relative to IPCC Sixth Assessment Report.
By region
Africa
In Africa, future population growth amplifies risks from sea level rise. Some 54.2 million people lived in the highly exposed low elevation coastal zones (LECZ) around 2000. This number will effectively double to around 110 million people by 2030, and then reach 185 to 230 million people by 2060. By then, the average regional sea level rise will be around 21 cm, with little difference from climate change scenarios. By 2100, Egypt, Mozambique and Tanzania are likely to have the largest number of people affected by annual flooding amongst all African countries. And under RCP8.5, 10 important cultural sites would be at risk of flooding and erosion by the end of the century.
In the near term, some of the largest displacement is projected to occur in the East Africa region. At least 750,000 people there are likely to be displaced from the coasts between 2020 and 2050. By 2050, 12 major African cities would collectively sustain cumulative damages of US$65 billion for the "moderate" climate change scenario RCP4.5 and between US$86.5 billion to US$137.5 billion on average: in the worst case, these damages could effectively triple. In all of these estimates, around half of the damages would occur in the Egyptian city of Alexandria. Hundreds of thousands of people in its low-lying areas may already need relocation in the coming decade. Across sub-Saharan Africa as a whole, damage from sea level rise could reach 2–4% of GDP by 2050, although this depends on the extent of future economic growth and climate change adaptation.
Asia
Asia has the largest population at risk from sea level due to its dense coastal populations. As of 2022, some 63 million people in East and South Asia were already at risk from a 100-year flood. This is largely due to inadequate coastal protection in many countries. Bangladesh, China, India, Indonesia, Japan, Pakistan, the Philippines, Thailand and Vietnam alone account for 70% of people exposed to sea level rise during the 21st century. Sea level rise in Bangladesh is likely to displace 0.9-2.1 million people by 2050. It may also force the relocation of up to one third of power plants as early as 2030, and many of the remaining plants would have to deal with the increased salinity of their cooling water. Nations like Bangladesh, Vietnam and China with extensive rice production on the coast are already seeing adverse impacts from saltwater intrusion.
Modelling results predict that Asia will suffer direct economic damages of US$167.6 billion at 0.47 meters of sea level rise. This rises to US$272.3 billion at 1.12 meters and US$338.1 billion at 1.75 meters. There is an additional indirect impact of US$8.5, 24 or 15 billion from population displacement at those levels. China, India, the Republic of Korea, Japan, Indonesia and Russia experience the largest economic losses. Out of the 20 coastal cities expected to see the highest flood losses by 2050, 13 are in Asia. Nine of these are the so-called sinking cities, where subsidence (typically caused by unsustainable groundwater extraction in the past) would compound sea level rise. These are Bangkok, Guangzhou, Ho Chi Minh City, Jakarta, Kolkata, Nagoya, Tianjin, Xiamen and Zhanjiang.
By 2050, Guangzhou would see 0.2 meters of sea level rise and estimated annual economic losses of US$254 million – the highest in the world. In Shanghai, coastal inundation amounts to about 0.03% of local GDP, yet would increase to 0.8% by 2100 even under the "moderate" RCP4.5 scenario in the absence of adaptation. The city of Jakarta is sinking so much (up to per year between 1982 and 2010 in some areas) that in 2019, the government had committed to relocate the capital of Indonesia to another city.
Australasia
In Australia, erosion and flooding of Queensland's Sunshine Coast beaches is likely to intensify by 60% by 2030. Without adaptation there would be a big impact on tourism. Adaptation costs for sea level rise would be three times higher under the high-emission RCP8.5 scenario than in the low-emission RCP2.6 scenario. Sea level rise of 0.2-0.3 meters is likely by 2050. In these conditions what is currently a 100-year flood would occur every year in the New Zealand cities of Wellington and Christchurch. With 0.5 m sea level rise, a current 100-year flood in Australia would occur several times a year. In New Zealand this would expose buildings with a collective worth of NZ$12.75 billion to new 100-year floods. A meter or so of sea level rise would threaten assets in New Zealand with a worth of NZD$25.5 billion. There would be a disproportionate impact on Maori-owned holdings and cultural heritage objects. Australian assets worth AUS$164–226 billion including many unsealed roads and railway lines would also be at risk. This amounts to a 111% rise in Australia's inundation costs between 2020 and 2100.
Central and South America
By 2100, coastal flooding and erosion will affect at least 3-4 million people in South America. Many people live in low-lying areas exposed to sea level rise. This includes 6% of the population of Venezuela, 56% of the population of Guyana and 68% of the population of Suriname. In Guyana much of the capital Georgetown is already below sea level. In Brazil, the coastal ecoregion of Caatinga is responsible for 99% of its shrimp production. A combination of sea level rise, ocean warming and ocean acidification threaten its unique ecosystem. Extreme wave or wind behavior disrupted the port complex of Santa Catarina 76 times in one 6-year period in the 2010s. There was a US$25,000-50,000 loss for each idle day. In Port of Santos, storm surges were three times more frequent between 2000 and 2016 than between 1928 and 1999.
Europe
Many sandy coastlines in Europe are vulnerable to erosion due to sea level rise. In Spain, Costa del Maresme is likely to retreat by 16 meters by 2050 relative to 2010. This could amount to 52 meters by 2100 under RCP8.5 Other vulnerable coastlines include the Tyrrhenian Sea coast of Italy's Calabria region, the Barra-Vagueira coast in Portugal and Nørlev Strand in Denmark.
In France, it was estimated that 8,000-10,000 people would be forced to migrate away from the coasts by 2080. The Italian city of Venice is located on islands. It is highly vulnerable to flooding and has already spent $6 billion on a barrier system. A quarter of the German state of Schleswig-Holstein, inhabited by over 350,000 people, is at low elevation and has been vulnerable to flooding since preindustrial times. Many levees already exist. Because of its complex geography, the authorities chose a flexible mix of hard and soft measures to cope with sea level rise of over 1 meter per century. In the United Kingdom, sea level at the end of the century would increase by 53 to 115 centimeters at the mouth of the River Thames and 30 to 90 centimeters at Edinburgh. The UK has divided its coast into 22 areas, each covered by a Shoreline Management Plan. Those are sub-divided into 2000 management units, working across three periods of 0–20, 20-50 and 50–100 years.
The Netherlands is a country that sits partially below sea level and is subsiding. It has responded by extending its Delta Works program. Drafted in 2008, the Delta Commission report said that the country must plan for a rise in the North Sea up to by 2100 and plan for a rise by 2200. It advised annual spending between €1.0 and €1.5 billion. This would support measures such as broadening coastal dunes and strengthening sea and river dikes. Worst-case evacuation plans were also drawn up.
North America
As of 2017, around 95 million Americans lived on the coast. The figures for Canada and Mexico were 6.5 million and 19 million. Increased chronic nuisance flooding and king tide flooding is already a problem in the highly vulnerable state of Florida. The US East Coast is also vulnerable. On average, the number of days with tidal flooding in the US increased 2 times in the years 2000–2020, reaching 3–7 days per year. In some areas the increase was much stronger: 4 times in the Southeast Atlantic and 11 times in the Western Gulf. By the year 2030 the average number is expected to be 7–15 days, reaching 25–75 days by 2050. U.S. coastal cities have responded with beach nourishment or beach replenishment. This trucks in mined sand in addition to other adaptation measures such as zoning, restrictions on state funding, and building code standards.
Along an estimated ~15% of the US coastline, the majority of local groundwater levels are already below sea level. This places those groundwater reservoirs at risk of sea water intrusion. That would render fresh water unusable once its concentration exceeds 2-3%. Damage is also widespread in Canada. It will affect major cities like Halifax and more remote locations like Lennox Island. The Mi'kmaq community there is already considering relocation due to widespread coastal erosion. In Mexico, damage from SLR to tourism hotspots like Cancun, Isla Mujeres, Playa del Carmen, Puerto Morelos and Cozumel could amount to US$1.4–2.3 billion. The increase in storm surge due to sea level rise is also a problem. Due to this effect Hurricane Sandy caused an additional US$8 billion in damage, impacted 36,000 more houses and 71,000 more people. In the future, the northern Gulf of Mexico, Atlantic Canada and the Pacific coast of Mexico would experience the greatest sea level rise. By 2030, flooding along the US Gulf Coast could cause economic losses of up to US$176 billion. Using nature-based solutions like wetland restoration and oyster reef restoration could avoid around US$50 billion of this.
By 2050, coastal flooding in the US is likely to rise tenfold to four "moderate" flooding events per year. That forecast is even without storms or heavy rainfall. In New York City, current 100-year flood would occur once in 19–68 years by 2050 and 4–60 years by 2080. By 2050, 20 million people in the greater New York City area would be at risk. This is because 40% of existing water treatment facilities would be compromised and 60% of power plants will need relocation.
By 2100, sea level rise of and would threaten 4.2 and 13.1 million people in the US, respectively. In California alone, of SLR could affect 600,000 people and threaten over US$150 billion in property with inundation. This potentially represents over 6% of the state's GDP. In North Carolina, a meter of SLR inundates 42% of the Albemarle-Pamlico Peninsula, costing up to US$14 billion. In nine southeast US states, the same level of sea level rise would claim up to 13,000 historical and archaeological sites, including over 1000 sites eligible for inclusion in the National Register for Historic Places.
Island nations
Small island states are nations with populations on atolls and other low islands. Atolls on average reach above sea level. These are the most vulnerable places to coastal erosion, flooding and salt intrusion into soils and freshwater caused by sea level rise. Sea level rise may make an island uninhabitable before it is completely flooded. Already, children in small island states encounter hampered access to food and water. They suffer an increased rate of mental and social disorders due to these stresses. At current rates, sea level rise would be high enough to make the Maldives uninhabitable by 2100. Five of the Solomon Islands have already disappeared due to the effects of sea level rise and stronger trade winds pushing water into the Western Pacific.
Adaptation to sea level rise is costly for small island nations as a large portion of their population lives in areas that are at risk. Nations like Maldives, Kiribati and Tuvalu already have to consider controlled international migration of their population in response to rising seas. The alternative of uncontrolled migration threatens to worsen the humanitarian crisis of climate refugees. In 2014, Kiribati purchased 20 square kilometers of land (about 2.5% of Kiribati's current area) on the Fijian island of Vanua Levu to relocate its population once their own islands are lost to the sea.
Fiji also suffers from sea level rise. It is in a comparatively safer position. Its residents continue to rely on local adaptation like moving further inland and increasing sediment supply to combat erosion instead of relocating entirely. Fiji has also issued a green bond of $50 million to invest in green initiatives and fund adaptation efforts. It is restoring coral reefs and mangroves to protect against flooding and erosion. It sees this as a more cost-efficient alternative to building sea walls. The nations of Palau and Tonga are taking similar steps. Even when an island is not threatened with complete disappearance from flooding, tourism and local economies may end up devastated. For instance, sea level rise of would cause partial or complete inundation of 29% of coastal resorts in the Caribbean. A further 49–60% of coastal resorts would be at risk from resulting coastal erosion.
| Physical sciences | Oceanography | Earth science |
21174200 | https://en.wikipedia.org/wiki/Gonggong%20%28dwarf%20planet%29 | Gonggong (dwarf planet) | Gonggong (minor-planet designation: 225088 Gonggong) is a dwarf planet and a member of the scattered disc beyond Neptune. It has a highly eccentric and inclined orbit during which it ranges from from the Sun. , its distance from the Sun is , and it is the sixth-farthest known Solar System object. According to the Deep Ecliptic Survey, Gonggong is in a 3:10 orbital resonance with Neptune, in which it completes three orbits around the Sun for every ten orbits completed by Neptune. Gonggong was discovered in July 2007 by American astronomers Megan Schwamb, Michael Brown, and David Rabinowitz at the Palomar Observatory, and the discovery was announced in January 2009.
At approximately in diameter, Gonggong is similar in size to Pluto's moon Charon, making it the fifth-largest known trans-Neptunian object (apart possibly from Charon). It may be sufficiently massive to be in hydrostatic equilibrium and therefore a dwarf planet. Gonggong's large mass makes retention of a tenuous atmosphere of methane just possible, though such an atmosphere would slowly escape into space. The object is named after Gònggōng, a Chinese water god responsible for chaos, floods and the tilt of the Earth. The name was chosen by its discoverers in 2019, when they hosted an online poll for the general public to help choose a name for the object, and the name Gonggong won.
Gonggong is red, likely due to the presence of organic compounds called tholins on its surface. Water ice is also present on its surface, which hints at a brief period of cryovolcanic activity in the distant past. With a rotation period of around 22 hours, Gonggong rotates slowly compared to other trans-Neptunian objects, which typically have periods of less than 12 hours. The slow rotation of Gonggong may have been caused by tidal forces from its natural satellite, named Xiangliu.
History
Discovery
Gonggong was discovered by American astronomers Megan Schwamb, Michael Brown and David Rabinowitz on 17 July 2007. The discovery was part of the Palomar Distant Solar System Survey, a survey conducted to find distant objects in the region of , beyond from the Sun, using the Samuel Oschin telescope at Palomar Observatory near San Diego, California. The survey was designed to detect the movements of objects out to at least 1,000 AU from the Sun. Schwamb identified Gonggong by comparing images using the blinking technique. In the discovery images, Gonggong appeared to move slowly, suggesting that it is a distant object. The discovery was part of Schwamb's doctoral thesis. At that time, Schwamb was a graduate student of Michael Brown at the California Institute of Technology.
Gonggong was formally announced in a Minor Planet Electronic Circular on 7 January 2009. It was then given the provisional designation because it was discovered during the second half of July 2007. The last letter and numbers of its designation indicate that it is the 267th object discovered during the latter half of July. , it has been observed 230 times over 13 oppositions, and has been identified in two precovery images, with the earliest image taken by the La Silla Observatory on 19 August 1985.
Name and symbol
The object is named after Gonggong, a water god in Chinese mythology. Gonggong is depicted as having a copper-and-iron, red-haired human head (or sometimes torso) and the body or tail of a serpent. Gonggong was responsible for creating chaos and catastrophe, causing flooding and tilting the Earth, until he was sent into exile. Gonggong is often accompanied by his minister, Xiangliu, a nine-headed poisonous snake monster who was also responsible for causing flooding and destruction.
Before its official naming, Gonggong was the largest known unnamed object in the Solar System. Initially after the discovery of Gonggong, Brown nicknamed the object "Snow White" for its presumed white color based on his assumption that it may be a member of the icy Haumea collisional family. The nickname also fit because, by that time, Brown's team had discovered seven other large trans-Neptunian objects which were collectively referred to as the "seven dwarfs": in 2002, in 2003, , and in 2004, and and in 2005. However, Gonggong turned out to be very red in color, comparable to Quaoar, so the nickname was dropped. On 2 November 2009, two years after its discovery, the Minor Planet Center assigned the minor planet number 225088 to Gonggong.
When Gonggong's discovery was first announced, Brown did not name it, as he considered it to be an unremarkable object, despite its large size. In 2011, he declared that he now had enough information to justify naming it, because of the discovery of water ice and the possibility of methane on its surface, which made it noteworthy enough to warrant further study. Following the Kepler spacecraft's large revision of Gonggong's size in 2016, Schwamb justified that Gonggong was eligible for naming, an acknowledgement of its large size and that its characteristics were known with enough certainty for a name to be given to reflect them.
In 2019, the discoverers of Gonggong hosted an online poll for the general public to choose between three possible names: Gonggong (Chinese), Holle (German), and Vili (Norse). These were selected by the discoverers in accordance with the International Astronomical Union's (IAU's) minor planet naming criteria, which state that objects with orbits like that of Gonggong must be given names related to mythological figures that are associated with creation. The three options were chosen because they were associated with water, ice, snow, and the color red—all characteristics of Gonggong—and because they had associated figures that could later provide a name for Gonggong's satellite. The name for Gonggong's satellite was not chosen by the hosts of the naming poll, as this privilege is reserved for its discoverers.
Having gained 46 percent of the 280,000 votes, on 29 May 2019, the discovery team announced Gonggong as the winning name. The name was proposed to the IAU's Committee on Small Body Nomenclature (CSBN), which is responsible for naming minor planets. The name was accepted by the CSBN and was announced by the Minor Planet Center on 5 February 2020.
As planetary symbols are no longer used regularly in astronomy, Gonggong never received a symbol in the astronomical literature. A symbol , used mostly among astrologers, is included in Unicode as . The symbol was designed by Denis Moskowitz, a software engineer in Massachusetts; it combines the Chinese character 共 gòng with a snake's tail.
Orbit
Gonggong orbits the Sun at an average distance of , and completes a full orbit in 554 years. The orbit of Gonggong is highly inclined to the ecliptic, with an orbital inclination of 30.7 degrees. Its orbit is also highly eccentric, with a measured orbital eccentricity of 0.50. Due to its highly eccentric orbit, the distance of Gonggong from the Sun varies greatly over the course of its orbit, from at aphelion, its furthest point from the Sun, to around at perihelion, its closest point to the Sun. Gonggong last reached perihelion in 1857, and is currently moving farther from the Sun, toward its aphelion. Gonggong will reach aphelion by 2134.
The period, inclination and eccentricity of Gonggong's orbit are all rather extreme compared to other large bodies in the Solar System. Among likely dwarf planets, its period is the third-longest, at 554 years compared to 558 years for and the ca. 11,400 years of . Its 31° inclination is second, after 44° for Eris, and its 0.50 eccentricity is also (a rather distant) second, after Sedna at 0.84.
The Minor Planet Center lists it as a scattered disc object for its eccentric and distant orbit. The Deep Ecliptic Survey shows the orbit of Gonggong to be in a 3:10 resonance with Neptune; Gonggong completes three orbits around the Sun for every ten orbits completed by Neptune.
, Gonggong is about from the Sun and is moving away at a speed of . It is the eleventh-farthest known Solar System object from the Sun, preceding (89.5 AU), (89.6 AU), (90.3 AU), (92.4 AU), (95.9 AU), (97.2 AU), (99.0 AU), (111.0 AU), (123.5 AU), and (~ 132 AU). Gonggong is more distant than , which is located 84.3 AU from the Sun . It has been farther from the Sun than Sedna since 2013, and it will surpass Eris in distance by 2045.
Brightness
Gonggong has an absolute magnitude (H) of 2.34, which makes it the seventh-brightest trans-Neptunian object known. It is dimmer than (H=2.31; D=917 km) but brighter than Quaoar (H=2.82; D=1,110 km). The Minor Planet Center and the Jet Propulsion Laboratory Small-Body Database assume a brighter absolute magnitude of 1.6 and 1.8, respectively, which would make it the fifth brightest trans-Neptunian object.
Being 88 AU from the Sun, the apparent magnitude of Gonggong is only 21.5, and so it is too dim to be seen from Earth with the naked eye. Although closer to the Sun than the dwarf planet Eris, Gonggong appears dimmer, as Eris has a higher albedo and an apparent magnitude of 18.8.
Physical characteristics
Surface and spectra
The surface of Gonggong has an albedo (reflectivity) of 0.14. The surface composition and spectrum of Gonggong is expected to be similar to that of , as both objects are red in color and display signs of water ice and possibly methane in their spectra. The reflectance spectrum of Gonggong was first measured in 2011 at near-infrared wavelengths, with the Folded port InfraRed Echellette (FIRE) spectrograph on the Magellan Baade Telescope at the Las Campanas Observatory in Chile. Gonggong's spectrum exhibits a strong red spectral slope along with broad absorption bands at wavelengths of 1.5 μm and 2 μm, meaning that Gonggong reflects more light at these wavelengths. Additional photometric measurements from the Hubble Space Telescope's Wide Field Camera 3 instrument display similar absorption bands at 1.5 μm, which are characteristic features of water ice, a substance often found on large Kuiper belt objects. The presence of water ice on the surface of Gonggong implies a brief period of cryovolcanism in the distant past, when water erupted from its interior, deposited onto its surface, and subsequently froze.
Gonggong is among the reddest trans-Neptunian objects known, especially in the visible and near-infrared. Its red color is unexpected for an object with a substantial amount of water ice on its surface, which are typically neutral in color, hence why Gonggong was initially nicknamed "Snow White". Gonggong's color implies that methane is present on its surface, although it was not directly detected in the spectrum of Gonggong due to the low signal-to-noise ratio of the data. The presence of methane frost would account for its color, as a result of the photolysis of methane by solar radiation and cosmic rays producing reddish organic compounds known as tholins. Observations of Gonggong's near-infrared spectrum in 2015 revealed an absorption feature at 2.27 μm, indicating the presence of methanol along with its irradiation products on its surface.
Gonggong is large enough to be able to retain trace amounts of volatile methane on its surface, even when at its closest distance to the Sun (33.7 AU), where temperatures are higher than that of Quaoar. In particular, the large size of Gonggong means that it is likely to retain trace amounts of other volatiles, including ammonia, carbon monoxide, and possibly nitrogen, which almost all trans-Neptunian objects lose over the course of their existence. Like Quaoar, Gonggong is expected to be near the mass limit at which it is able to retain those volatile materials on its surface.
In 2022, low resolution near-infrared (0.7–5 μm) spectroscopic observations by the James Webb Space Telescope (JWST) revealed the presence of significant amounts of ethane ice (C2H6) on the surface of Gonggong, though there appears to be less ethane on Gonggong than on Sedna. The JWST spectra also contain evidence of presence of small amounts of carbon dioxide (CO2) complexed with either dark surface material or some ices as well as complex organics. On the other hand no evidence of presence of methane (CH4) and methanol (CH3OH) was found at variance with the earlier observations.
Atmosphere
The presence of tholins on the surface of Gonggong implies the possible existence of a tenuous methane atmosphere, analogous to Quaoar. Although Gonggong occasionally comes closer to the Sun than Quaoar, where it becomes warm enough that a methane atmosphere should evaporate, its larger mass could make the retention of methane just possible. During aphelion, methane along with other volatiles would condense on Gonggong's surface, allowing for long-term irradiation that would otherwise result in a decrease in surface albedo. The lower surface albedo would contribute to the loss of highly volatile materials such as nitrogen, as a lower albedo corresponds to more light being absorbed by the surface rather than being reflected, thus resulting in greater surface heating. Hence, the nitrogen content of Gonggong's atmosphere is expected to be depleted to trace amounts while methane is likely retained.
Gonggong is thought to have had cryovolcanic activity along with a more substantial atmosphere shortly after its formation. Such cryovolcanic activity is expected to have been brief, and the resulting atmosphere gradually escaped over time. Volatile gases, such as nitrogen and carbon monoxide, were lost, while less volatile gases such as methane are likely to remain in its present tenuous atmosphere.
Size
As of 2019, Gonggong is estimated to have a diameter of , derived from radiometric measurements, its calculated mass, and assuming a density similar to other similar bodies. This would make Gonggong the fifth-largest trans-Neptunian object, after Pluto, Eris, Haumea and Makemake. Gonggong is approximately the size of Pluto's moon Charon, although Gonggong's current size estimate has an uncertainty of .
The International Astronomical Union (IAU) has not addressed the possibility of officially accepting additional dwarf planets since the acceptance of Makemake and Haumea in 2008, prior to the announcement of Gonggong in 2009. Gonggong is large enough to be considered a dwarf planet by several astronomers. Brown states that Gonggong "must be a dwarf planet even if predominantly rocky", based on the 2013 radiometric measurement of . Scott Sheppard and colleagues think that it is likely to be a dwarf planet, based on its minimum possible diameter— under the assumption of a completely reflective surface with an albedo of 1—and what was at the time the expected lower size limit of around for hydrostatic equilibrium in cold icy-rocky bodies. However, Iapetus is not in equilibrium despite being in diameter, so this remains just a possibility.
In 2010, astronomer Gonzalo Tancredi initially estimated Gonggong to have a very large diameter of , though its dwarf planet status was unclear as there was no lightcurve data or other information to ascertain its size. Gonggong is too distant to be resolved directly; Brown placed a rough estimate of its diameter ranging from , based on an albedo of 0.18 which was the best fit in his model. A survey led by a team of astronomers using the European Space Agency's Herschel Space Observatory in 2012 determined its diameter to be (), based on the thermal properties of Gonggong observed in the far infrared range. This measurement is consistent with Brown's estimate. Later observations in 2013 using combined thermal emission data from Herschel and the Spitzer Space Telescope suggested a smaller size of (), though this estimate had a larger range of uncertainty.
In 2016, combined observations from the Kepler spacecraft and archival thermal emission data from Herschel suggested that Gonggong was much larger than previously thought, giving a size estimate of () based on an assumed equator-on view and a lower estimated albedo of 0.089. This would have made Gonggong the third-largest trans-Neptunian object after Eris and Pluto, larger than Makemake (). These observations of Gonggong were part of the Kepler spacecraft's K2 mission which includes studying small Solar System bodies. Subsequent measurements in 2018 revised the size of Gonggong to (), based on the mass and density of Gonggong derived from the orbit of its satellite and the discovery that the viewing direction was almost pole-on. With this size estimate, Gonggong is again thought to be the fifth-largest trans-Neptunian object.
Mass, density and rotation
Based on the orbit of its satellite, the mass of Gonggong has been calculated to be , with a density of . Given the mass, the 2016 size estimate of would have implied an unexpectedly low (and likely erroneous) density of .
Gonggong is the fifth most massive trans-Neptunian object, after Eris, Pluto, Haumea, and Makemake. It is slightly more massive and denser than Charon, which has a mass of and a density of . Due to its large size, mass, and density, Gonggong is expected to be in hydrostatic equilibrium, taking the shape of a MacLaurin spheroid that is slightly flattened due to its rotation.
The rotation period of Gonggong was first measured in March 2016, through observations of variations in its brightness with the Kepler space telescope. Gonggong's light curve amplitude as observed by Kepler is small, only varying in brightness by about 0.09 magnitudes. The small light curve amplitude of Gonggong indicates that it is being viewed at a pole-on configuration, further evidenced by the observed inclined orbit of its satellite. The Kepler observations provided ambiguous values of and hours for the rotation period. Based on a best-fit model for its rotation pole orientation, the value of hours is thought to be the more plausible one. Gonggong rotates slowly compared to other trans-Neptunian objects, which usually have periods between 6 and 12 hours. Due to its slow rotation, it is expected to have a low oblateness of 0.03 or 0.007, for rotation periods of 22.4 or 44.81 hours, respectively.
Satellite
Following the March 2016 discovery that Gonggong was an unusually slow rotator, the possibility was raised that a satellite may have slowed it down via tidal forces. The indications of a possible satellite orbiting Gonggong led Csaba Kiss and his team to analyze archival Hubble observations of Gonggong. Their analysis of Hubble images taken on 18 September 2010 revealed a faint satellite orbiting Gonggong at a distance of at least . The discovery was announced in a Division for Planetary Sciences meeting on 17 October 2016. The satellite is approximately in diameter and has an orbital period of 25 days. On 5 February 2020 the satellite was officially named Xiangliu, after the nine-headed poisonous snake monster that accompanied Gonggong in Chinese mythology. This naming came at the same time that Gonggong itself was officially named.
Exploration
It was calculated by planetary scientist Amanda Zangari that a flyby mission to Gonggong would take a minimum of over 20 years with current rocket capabilities. A flyby mission could take just under 25 years using a Jupiter gravity assist, based on a launch date of 2030 or 2031. Gonggong would be approximately 95 AU from the Sun when the spacecraft arrives.
| Physical sciences | Solar System | Astronomy |
5010838 | https://en.wikipedia.org/wiki/Sine%20and%20cosine | Sine and cosine | In mathematics, sine and cosine are trigonometric functions of an angle. The sine and cosine of an acute angle are defined in the context of a right triangle: for the specified angle, its sine is the ratio of the length of the side that is opposite that angle to the length of the longest side of the triangle (the hypotenuse), and the cosine is the ratio of the length of the adjacent leg to that of the hypotenuse. For an angle , the sine and cosine functions are denoted as and .
The definitions of sine and cosine have been extended to any real value in terms of the lengths of certain line segments in a unit circle. More modern definitions express the sine and cosine as infinite series, or as the solutions of certain differential equations, allowing their extension to arbitrary positive and negative values and even to complex numbers.
The sine and cosine functions are commonly used to model periodic phenomena such as sound and light waves, the position and velocity of harmonic oscillators, sunlight intensity and day length, and average temperature variations throughout the year. They can be traced to the and functions used in Indian astronomy during the Gupta period.
Elementary descriptions
Right-angled triangle definition
To define the sine and cosine of an acute angle , start with a right triangle that contains an angle of measure ; in the accompanying figure, angle in a right triangle is the angle of interest. The three sides of the triangle are named as follows:
The opposite side is the side opposite to the angle of interest; in this case, it is .
The hypotenuse is the side opposite the right angle; in this case, it is . The hypotenuse is always the longest side of a right-angled triangle.
The adjacent side is the remaining side; in this case, it is . It forms a side of (and is adjacent to) both the angle of interest and the right angle.
Once such a triangle is chosen, the sine of the angle is equal to the length of the opposite side divided by the length of the hypotenuse, and the cosine of the angle is equal to the length of the adjacent side divided by the length of the hypotenuse:
The other trigonometric functions of the angle can be defined similarly; for example, the tangent is the ratio between the opposite and adjacent sides or equivalently the ratio between the sine and cosine functions. The reciprocal of sine is cosecant, which gives the ratio of the hypotenuse length to the length of the opposite side. Similarly, the reciprocal of cosine is secant, which gives the ratio of the hypotenuse length to that of the adjacent side. The cotangent function is the ratio between the adjacent and opposite sides, a reciprocal of a tangent function. These functions can be formulated as:
Special angle measures
As stated, the values and appear to depend on the choice of a right triangle containing an angle of measure . However, this is not the case as all such triangles are similar, and so the ratios are the same for each of them. For example, each leg of the 45-45-90 right triangle is 1 unit, and its hypotenuse is ; therefore, . The following table shows the special value of each input for both sine and cosine with the domain between . The input in this table provides various unit systems such as degree, radian, and so on. The angles other than those five can be obtained by using a calculator.
Laws
The law of sines is useful for computing the lengths of the unknown sides in a triangle if two angles and one side are known. Given that a triangle with sides , , and , and angles opposite those sides , , and . The law states,
This is equivalent to the equality of the first three expressions below:
where is the triangle's circumradius.
The law of cosines is useful for computing the length of an unknown side if two other sides and an angle are known. The law states,
In the case where from which , the resulting equation becomes the Pythagorean theorem.
Vector definition
The cross product and dot product are operations on two vectors in Euclidean vector space. The sine and cosine functions can be defined in terms of the cross product and dot product. If and are vectors, and is the angle between and , then sine and cosine can be defined as:
Analytic descriptions
Unit circle definition
The sine and cosine functions may also be defined in a more general way by using unit circle, a circle of radius one centered at the origin , formulated as the equation of in the Cartesian coordinate system. Let a line through the origin intersect the unit circle, making an angle of with the positive half of the axis. The and coordinates of this point of intersection are equal to and , respectively; that is,
This definition is consistent with the right-angled triangle definition of sine and cosine when because the length of the hypotenuse of the unit circle is always 1; mathematically speaking, the sine of an angle equals the opposite side of the triangle, which is simply the coordinate. A similar argument can be made for the cosine function to show that the cosine of an angle when , even under the new definition using the unit circle.
Graph of a function and its elementary properties
Using the unit circle definition has the advantage of drawing a graph of sine and cosine functions. This can be done by rotating counterclockwise a point along the circumference of a circle, depending on the input . In a sine function, if the input is , the point is rotated counterclockwise and stopped exactly on the axis. If , the point is at the circle's halfway. If , the point returned to its origin. This results that both sine and cosine functions have the range between .
Extending the angle to any real domain, the point rotated counterclockwise continuously. This can be done similarly for the cosine function as well, although the point is rotated initially from the coordinate. In other words, both sine and cosine functions are periodic, meaning any angle added by the circumference's circle is the angle itself. Mathematically,
A function is said to be odd if , and is said to be even if . The sine function is odd, whereas the cosine function is even. Both sine and cosine functions are similar, with their difference being shifted by . This means,
Zero is the only real fixed point of the sine function; in other words the only intersection of the sine function and the identity function is . The only real fixed point of the cosine function is called the Dottie number. The Dottie number is the unique real root of the equation . The decimal expansion of the Dottie number is approximately 0.739085.
Continuity and differentiation
The sine and cosine functions are infinitely differentiable. The derivative of sine is cosine, and the derivative of cosine is negative sine:
Continuing the process in higher-order derivative results in the repeated same functions; the fourth derivative of a sine is the sine itself. These derivatives can be applied to the first derivative test, according to which the monotonicity of a function can be defined as the inequality of function's first derivative greater or less than equal to zero. It can also be applied to second derivative test, according to which the concavity of a function can be defined by applying the inequality of the function's second derivative greater or less than equal to zero. The following table shows that both sine and cosine functions have concavity and monotonicity—the positive sign () denotes a graph is increasing (going upward) and the negative sign () is decreasing (going downward)—in certain intervals. This information can be represented as a Cartesian coordinates system divided into four quadrants.
Both sine and cosine functions can be defined by using differential equations. The pair of is the solution to the two-dimensional system of differential equations and with the initial conditions and . One could interpret the unit circle in the above definitions as defining the phase space trajectory of the differential equation with the given initial conditions. It can be interpreted as a phase space trajectory of the system of differential equations and starting from the initial conditions and .
Integral and the usage in mensuration
Their area under a curve can be obtained by using the integral with a certain bounded interval. Their antiderivatives are:
where denotes the constant of integration. These antiderivatives may be applied to compute the mensuration properties of both sine and cosine functions' curves with a given interval. For example, the arc length of the sine curve between and is
where is the incomplete elliptic integral of the second kind with modulus . It cannot be expressed using elementary functions. In the case of a full period, its arc length is
where is the gamma function and is the lemniscate constant.
Inverse functions
The inverse function of sine is arcsine or inverse sine, denoted as "arcsin", "asin", or . The inverse function of cosine is arccosine, denoted as "arccos", "acos", or . As sine and cosine are not injective, their inverses are not exact inverse functions, but partial inverse functions. For example, , but also , , and so on. It follows that the arcsine function is multivalued: , but also , , and so on. When only one value is desired, the function may be restricted to its principal branch. With this restriction, for each in the domain, the expression will evaluate only to a single value, called its principal value. The standard range of principal values for arcsin is from to , and the standard range for arccos is from to .
The inverse function of both sine and cosine are defined as:
where for some integer ,
By definition, both functions satisfy the equations:
and
Other identities
According to Pythagorean theorem, the squared hypotenuse is the sum of two squared legs of a right triangle. Dividing the formula on both sides with squared hypotenuse resulting in the Pythagorean trigonometric identity, the sum of a squared sine and a squared cosine equals 1:
Sine and cosine satisfy the following double-angle formulas:
The cosine double angle formula implies that sin2 and cos2 are, themselves, shifted and scaled sine waves. Specifically,
The graph shows both sine and sine squared functions, with the sine in blue and the sine squared in red. Both graphs have the same shape but with different ranges of values and different periods. Sine squared has only positive values, but twice the number of periods.
Series and polynomials
Both sine and cosine functions can be defined by using a Taylor series, a power series involving the higher-order derivatives. As mentioned in , the derivative of sine is cosine and that the derivative of cosine is the negative of sine. This means the successive derivatives of are , , , , continuing to repeat those four functions. The th derivative, evaluated at the point 0:
where the superscript represents repeated differentiation. This implies the following Taylor series expansion at . One can then use the theory of Taylor series to show that the following identities hold for all real numbers —where is the angle in radians. More generally, for all complex numbers:
Taking the derivative of each term gives the Taylor series for cosine:
Both sine and cosine functions with multiple angles may appear as their linear combination, resulting in a polynomial. Such a polynomial is known as the trigonometric polynomial. The trigonometric polynomial's ample applications may be acquired in its interpolation, and its extension of a periodic function known as the Fourier series. Let and be any coefficients, then the trigonometric polynomial of a degree —denoted as —is defined as:
The trigonometric series can be defined similarly analogous to the trigonometric polynomial, its infinite inversion. Let and be any coefficients, then the trigonometric series can be defined as:
In the case of a Fourier series with a given integrable function , the coefficients of a trigonometric series are:
Complex numbers relationship
Complex exponential function definitions
Both sine and cosine can be extended further via complex number, a set of numbers composed of both real and imaginary numbers. For real number , the definition of both sine and cosine functions can be extended in a complex plane in terms of an exponential function as follows:
Alternatively, both functions can be defined in terms of Euler's formula:
When plotted on the complex plane, the function for real values of traces out the unit circle in the complex plane. Both sine and cosine functions may be simplified to the imaginary and real parts of as:
When for real values and , where , both sine and cosine functions can be expressed in terms of real sines, cosines, and hyperbolic functions as:
Polar coordinates
Sine and cosine are used to connect the real and imaginary parts of a complex number with its polar coordinates :
and the real and imaginary parts are
where and represent the magnitude and angle of the complex number .
For any real number , Euler's formula in terms of polar coordinates is stated as .
Complex arguments
Applying the series definition of the sine and cosine to a complex argument, z, gives:
where sinh and cosh are the hyperbolic sine and cosine. These are entire functions.
It is also sometimes useful to express the complex sine and cosine functions in terms of the real and imaginary parts of its argument:
Partial fraction and product expansions of complex sine
Using the partial fraction expansion technique in complex analysis, one can find that the infinite series
both converge and are equal to . Similarly, one can show that
Using product expansion technique, one can derive
Usage of complex sine
sin(z) is found in the functional equation for the Gamma function,
which in turn is found in the functional equation for the Riemann zeta-function,
As a holomorphic function, sin z is a 2D solution of Laplace's equation:
The complex sine function is also related to the level curves of pendulums.
Complex graphs
Background
Etymology
The word sine is derived, indirectly, from the Sanskrit word 'bow-string' or more specifically its synonym (both adopted from Ancient Greek 'string'), due to visual similarity between the arc of a circle with its corresponding chord and a bow with its string (see jyā, koti-jyā and utkrama-jyā). This was transliterated in Arabic as , which is meaningless in that language and written as (). Since Arabic is written without short vowels, was interpreted as the homograph (جيب), which means 'bosom', 'pocket', or 'fold'. When the Arabic texts of Al-Battani and al-Khwārizmī were translated into Medieval Latin in the 12th century by Gerard of Cremona, he used the Latin equivalent sinus (which also means 'bay' or 'fold', and more specifically 'the hanging fold of a toga over the breast'). Gerard was probably not the first scholar to use this translation; Robert of Chester appears to have preceded him and there is evidence of even earlier usage. The English form sine was introduced in the 1590s.
The word cosine derives from an abbreviation of the Latin 'sine of the complementary angle' as cosinus in Edmund Gunter's Canon triangulorum (1620), which also includes a similar definition of cotangens.
History
While the early study of trigonometry can be traced to antiquity, the trigonometric functions as they are in use today were developed in the medieval period. The chord function was discovered by Hipparchus of Nicaea (180–125 BCE) and Ptolemy of Roman Egypt (90–165 CE).
The sine and cosine functions can be traced to the and functions used in Indian astronomy during the Gupta period (Aryabhatiya and Surya Siddhanta), via translation from Sanskrit to Arabic and then from Arabic to Latin.
All six trigonometric functions in current use were known in Islamic mathematics by the 9th century, as was the law of sines, used in solving triangles. With the exception of the sine (which was adopted from Indian mathematics), the other five modern trigonometric functions were discovered by Arabic mathematicians, including the cosine, tangent, cotangent, secant and cosecant. Al-Khwārizmī (c. 780–850) produced tables of sines, cosines and tangents. Muhammad ibn Jābir al-Harrānī al-Battānī (853–929) discovered the reciprocal functions of secant and cosecant, and produced the first table of cosecants for each degree from 1° to 90°.
The first published use of the abbreviations sin, cos, and tan is by the 16th-century French mathematician Albert Girard; these were further promulgated by Euler (see below). The Opus palatinum de triangulis of Georg Joachim Rheticus, a student of Copernicus, was probably the first in Europe to define trigonometric functions directly in terms of right triangles instead of circles, with tables for all six trigonometric functions; this work was finished by Rheticus' student Valentin Otho in 1596.
In a paper published in 1682, Leibniz proved that sin x is not an algebraic function of x. Roger Cotes computed the derivative of sine in his Harmonia Mensurarum (1722). Leonhard Euler's Introductio in analysin infinitorum (1748) was mostly responsible for establishing the analytic treatment of trigonometric functions in Europe, also defining them as infinite series and presenting "Euler's formula", as well as the near-modern abbreviations sin., cos., tang., cot., sec., and cosec.
Software implementations
There is no standard algorithm for calculating sine and cosine. IEEE 754, the most widely used standard for the specification of reliable floating-point computation, does not address calculating trigonometric functions such as sine. The reason is that no efficient algorithm is known for computing sine and cosine with a specified accuracy, especially for large inputs.
Algorithms for calculating sine may be balanced for such constraints as speed, accuracy, portability, or range of input values accepted. This can lead to different results for different algorithms, especially for special circumstances such as very large inputs, e.g. sin(10).
A common programming optimization, used especially in 3D graphics, is to pre-calculate a table of sine values, for example one value per degree, then for values in-between pick the closest pre-calculated value, or linearly interpolate between the 2 closest values to approximate it. This allows results to be looked up from a table rather than being calculated in real time. With modern CPU architectures this method may offer no advantage.
The CORDIC algorithm is commonly used in scientific calculators.
The sine and cosine functions, along with other trigonometric functions, are widely available across programming languages and platforms. In computing, they are typically abbreviated to sin and cos.
Some CPU architectures have a built-in instruction for sine, including the Intel x87 FPUs since the 80387.
In programming languages, sin and cos are typically either a built-in function or found within the language's standard math library. For example, the C standard library defines sine functions within math.h: sin(double), sinf(float), and sinl(long double). The parameter of each is a floating point value, specifying the angle in radians. Each function returns the same data type as it accepts. Many other trigonometric functions are also defined in math.h, such as for cosine, arc sine, and hyperbolic sine (sinh). Similarly, Python defines math.sin(x) and math.cos(x) within the built-in math module. Complex sine and cosine functions are also available within the cmath module, e.g. cmath.sin(z). CPython's math functions call the C math library, and use a double-precision floating-point format.
Turns based implementations
Some software libraries provide implementations of sine and cosine using the input angle in half-turns, a half-turn being an angle of 180 degrees or radians. Representing angles in turns or half-turns has accuracy advantages and efficiency advantages in some cases. In MATLAB, OpenCL, R, Julia, CUDA, and ARM, these functions are called sinpi and cospi. For example, sinpi(x) would evaluate to where x is expressed in half-turns, and consequently the final input to the function, can be interpreted in radians by .
The accuracy advantage stems from the ability to perfectly represent key angles like full-turn, half-turn, and quarter-turn losslessly in binary floating-point or fixed-point. In contrast, representing , , and in binary floating-point or binary scaled fixed-point always involves a loss of accuracy since irrational numbers cannot be represented with finitely many binary digits.
Turns also have an accuracy advantage and efficiency advantage for computing modulo to one period. Computing modulo 1 turn or modulo 2 half-turns can be losslessly and efficiently computed in both floating-point and fixed-point. For example, computing modulo 1 or modulo 2 for a binary point scaled fixed-point value requires only a bit shift or bitwise AND operation. In contrast, computing modulo involves inaccuracies in representing .
For applications involving angle sensors, the sensor typically provides angle measurements in a form directly compatible with turns or half-turns. For example, an angle sensor may count from 0 to 4096 over one complete revolution. If half-turns are used as the unit for angle, then the value provided by the sensor directly and losslessly maps to a fixed-point data type with 11 bits to the right of the binary point. In contrast, if radians are used as the unit for storing the angle, then the inaccuracies and cost of multiplying the raw sensor integer by an approximation to would be incurred.
| Mathematics | Specific functions | null |
5012175 | https://en.wikipedia.org/wiki/Albatross | Albatross | Albatrosses, of the biological family Diomedeidae, are large seabirds related to the procellariids, storm petrels, and diving petrels in the order Procellariiformes (the tubenoses). They range widely in the Southern Ocean and the North Pacific. They are absent from the North Atlantic, although fossil remains of short-tailed albatross show they once lived there up to the Pleistocene, and occasional vagrants are found. Great albatrosses are among the largest of flying birds, with wingspans reaching up to and bodies over in length. The albatrosses are usually regarded as falling into four genera, but disagreement exists over the number of species.
Albatrosses are highly efficient in the air, using dynamic soaring and slope soaring to cover great distances with little exertion. They feed on squid, fish, and krill by either scavenging, surface seizing, or diving. Albatrosses are colonial, nesting for the most part on remote oceanic islands, often with several species nesting together. Pair bonds between males and females form over several years, with the use of "ritualised dances", and last for the life of the pair. A breeding season can take over a year from laying to fledging, with a single egg laid in each breeding attempt. A Laysan albatross, named Wisdom, on Midway Island is the oldest-known wild bird in the world; she was first banded in 1956 by Chandler Robbins.
Of the 22 species of albatrosses recognised by the IUCN, 21 are listed as at some level of concern; two species are Critically Endangered, seven species are Endangered, six species are Vulnerable, and six species are Near Threatened. Numbers of albatrosses have declined in the past due to harvesting for feathers. Albatrosses are threatened by introduced species, such as rats and feral cats that attack eggs, chicks, and nesting adults; by pollution; by a serious decline in fish stocks in many regions largely due to overfishing; and by longline fishing. Longline fisheries pose the greatest threat, as feeding birds are attracted to the bait, become hooked on the lines, and drown. Identified stakeholders such as governments, conservation organisations, and people in the fishing industry are all working toward reducing this phenomenon.
Etymology
The name "Albatross" is derived from the Arabic al-qādūs or al-ḡaṭṭās (a pelican; literally, "the diver"), which travelled to English via the Portuguese form Alcatraz ("gannet"), which is also the origin of the name of the former prison Alcatraz. The Oxford English Dictionary notes that the word Alcatraz was originally applied to the frigatebird; the modification to albatross was perhaps influenced by Latin Albus, meaning "white", in contrast to frigatebirds, which are black.
They were once commonly known as goonie birds or gooney birds, particularly those of the North Pacific. In the Southern Hemisphere, the name mollymawk is still well established in some areas, which is a corrupted form of malle-mugge, an old Dutch name for the northern fulmar. The name Diomedea, assigned to the albatrosses by Linnaeus, references the mythical metamorphosis of the companions of the Greek warrior Diomedes into birds. Finally, the name for the order, Procellariiformes, comes from the Latin word procella meaning "a violent wind" or "a storm".
Taxonomy and evolution
The "albatross" designation comprises between 13 and 24 species (the number is still a matter of some debate, with 21 being the most commonly accepted number) in four genera. These genera are the great albatrosses (Diomedea), the mollymawks (Thalassarche), the North Pacific albatrosses (Phoebastria), and the sooty albatrosses or sooties (Phoebetria). The North Pacific albatrosses are considered to be a sister taxon to the great albatrosses, while the sooty albatrosses are considered closer to the mollymawks.
The taxonomy of the albatross group has been a source of much debate. The Sibley-Ahlquist taxonomy places seabirds, birds of prey, and many others in a greatly enlarged order, the Ciconiiformes, whereas the ornithological organisations in North America, Europe, South Africa, Australia, and New Zealand retain the more traditional order Procellariiformes. The albatrosses can be separated from the other Procellariiformes both genetically and through morphological characteristics, size, their legs, and the arrangement of their nasal tubes (see below: Morphology and flight).
Within the family, the assignment of genera has been debated for over 100 years. Originally placed into a single genus, Diomedea, they were rearranged by Reichenbach into four different genera in 1852, then lumped back together and split apart again several times, acquiring 12 different genus names in total (though never more than eight at one time) by 1965 (Diomedea, Phoebastria, Thalassarche, Phoebetria, Thalassageron, Diomedella, Nealbatrus, Rhothonia, Julietata, Galapagornis, Laysanornis, and Penthirenia).
By 1965, in an attempt to bring some order back to the classification of albatrosses, they were lumped into two genera, Phoebetria (the sooty albatrosses, which most closely seemed to resemble the procellarids and were at the time considered "primitive" ) and Diomedea (the rest). Though a case was made for the simplification of the family (particularly the nomenclature), the classification was based on the morphological analysis by Elliott Coues in 1866, and paid little attention to more recent studies and even ignored some of Coues's suggestions.
Research by Gary Nunn of the American Museum of Natural History (1996) and other researchers around the world studied the mitochondrial DNA of all 14 accepted species, finding four, not two, monophyletic groups within the albatrosses. They proposed the resurrection of two of the old genus names, Phoebastria for the North Pacific albatrosses and Thalassarche for the mollymawks, with the great albatrosses retaining Diomedea and the sooty albatrosses staying in Phoebetria.
While some agree on the number of genera, fewer agree on the number of species. Historically, up to 80 different taxa have been described by different researchers; most of these were incorrectly identified juvenile birds.
Based on the work on albatross genera, Robertson and Nunn went on in 1998 to propose a revised taxonomy with 24 different species, compared to the 14 then accepted. This expanded taxonomy elevated many established subspecies to full species, but was criticised for not using, in every case, peer reviewed information to justify the splits. Since then, further studies have in some instances supported or disproved the splits; a 2004 paper analysing the mitochondrial DNA and microsatellites agreed with the conclusion that the Antipodean albatross and the Tristan albatross were distinct from the wandering albatross, per Robertson and Nunn, but found that the suggested Gibson's albatross, Diomedea gibsoni, was not distinct from the Antipodean albatross. For the most part, an interim taxonomy of 21 species is accepted by ITIS and many other researchers, though by no means all—in 2004 Penhallurick and Wink called for the number of species to be reduced to 13 (including the lumping of the Amsterdam albatross with the wandering albatross), although this paper was itself controversial.
Sibley and Ahlquist's molecular study of the evolution of the bird families has put the radiation of the Procellariiformes in the Oligocene period 35–30 million years ago (Mya), though this group probably originated earlier, with a fossil sometimes attributed to the order, a seabird known as Tytthostonyx, being found in late Cretaceous rocks (70 Mya). The molecular evidence suggests that the storm petrels were the first to diverge from the ancestral stock, and the albatrosses next, with the procellarids and diving petrels separating later. The earliest fossil albatrosses were found in Eocene to Oligocene rocks, although some of these are only tentatively assigned to the family and none appear to be particularly close to the living forms. They are Murunkus (Middle Eocene of Uzbekistan), Manu (early Oligocene of New Zealand), and an undescribed form from the Late Oligocene of South Carolina. The oldest widely accepted fossil albatross is Tydea septentrionalis from the early Oligocene of Belgium. Diomedavus knapptonensis is smaller than all extant albatrosses and was found in late Oligocene strata of Washington State, USA. Plotornis was formerly often considered a petrel but is now accepted as an albatross. It is from the Middle Miocene of France, a time when the split between the four modern genera was already underway as evidenced by Phoebastria californica and Diomedea milleri, both being mid-Miocene species from Sharktooth Hill, California. These show that the split between the great albatrosses and the North Pacific albatrosses occurred by 15 Mya. Similar fossil finds in the Southern Hemisphere put the split between the sooties and mollymawks at 10 Mya.
The fossil record of the albatrosses in the Northern Hemisphere is more complete than that of the Southern, and many fossil forms of albatross have been found in the North Atlantic, which today has no albatrosses. The remains of a colony of short-tailed albatrosses have been uncovered on the island of Bermuda, and the majority of fossil albatrosses from the North Atlantic have been of the genus Phoebastria (the North Pacific albatrosses); one, Phoebastria anglica, has been found in deposits in both North Carolina and England. Due to convergent evolution in particular of the leg and foot bones, remains of the prehistoric pseudotooth birds (Pelagornithidae) may be mistaken for those of extinct albatrosses; Manu may be such a case, and quite certainly the supposed giant albatross femur from the Early Pleistocene Dainichi Formation at Kakegawa, Japan, actually is from one of the last pseudotooth birds. Aldiomedes angustirostris was a uniquely narrow-beaked species from the Pliocene of New Zealand.
Morphology and flight
The albatrosses are a group of large to very large birds; they are the largest of the Procellariiformes. The bill is large, strong, and sharp-edged, with the upper mandible terminating in a large hook. This bill is composed of several horny plates, and along the sides are the two "tubes", long nostrils that give the order its former name (Tubinares, or tubenoses). The tubes of all albatrosses are along the sides of the bill, unlike the rest of the Procellariiformes, where the tubes run along the top of the bill. These tubes allow the albatrosses to measure the exact airspeed in flight; the nostrils are analogous to the pitot tubes in modern aircraft. The albatross needs accurate airspeed measurement to perform dynamic soaring. Like other Procellariiformes, they use their uniquely developed sense of smell to locate potential food sources, whereas most birds depend on eyesight. The feet have no hind toe and the three anterior toes are completely webbed. The legs are strong for the Procellariiformes, making them and the giant petrels the only members of that order that can walk well on land.
Albatrosses, along with all Procellariiformes, must excrete the salts they ingest in drinking sea water and eating marine invertebrates. All birds have an enlarged nasal gland at the base of the bill, above their eyes. This gland is inactive in species that do not require it, but in the Procellariiformes, it acts as a salt gland. Scientists are uncertain as to its exact processes, but do know in general terms that it removes salt by secreting a 5% saline solution that drips out of their noses or is forcibly ejected.
The adult plumage of most of the albatrosses is usually some variation of dark upper-wing and back with white undersides, often compared to that of a gull. The extent of colouration varies: the southern royal albatross is almost completely white except for the ends and trailing edges of the wings in fully mature males, while the Amsterdam albatross has an almost juvenile-like breeding plumage with a great deal of brown, particularly a strong brown band around the chest. Several species of mollymawks and North Pacific albatrosses have face markings like eye patches or have grey or yellow on the head and nape. Three albatross species, the black-footed albatross and the two sooty albatrosses, vary completely from the usual patterns and are almost entirely dark brown (or dark grey in places in the case of the light-mantled albatross). Albatrosses take several years to get their full adult breeding plumage.
The wingspans of the largest great albatrosses (genus Diomedea) are the largest of any bird, exceeding , although the other species' wingspans are considerably smaller, at as low as . The wings are stiff and cambered, with thickened, streamlined leading edges. Albatrosses travel long distances with two techniques used by many long-winged seabirds – dynamic soaring and slope soaring. Dynamic soaring involves repeatedly rising into wind and descending downwind, thus gaining energy from the vertical wind gradient. The only effort expended is in the turns at the top and bottom of every such loop. This maneuver allows the bird to cover almost without flapping its wings. Slope soaring uses the rising air on the windward side of large waves. Albatross have high glide ratios, around 22:1 to 23:1, meaning that for every metre they drop, they can travel forward twenty-two metres. They are aided in soaring by a shoulder-lock, a sheet of tendon that locks the wing when fully extended, allowing the wing to be kept outstretched without any muscle expenditure, a morphological adaptation they share with the giant petrels.
Albatrosses combine these soaring techniques with the use of predictable weather systems; albatrosses in the Southern Hemisphere flying north from their colonies take a clockwise route, and those flying south fly counterclockwise. Albatrosses are so well adapted to this lifestyle that their heart rates while flying are close to their basal heart rate when resting. This efficiency is such that the most energetically demanding aspect of a foraging trip is not the distance covered, but the landings, take-offs and hunting they undertake having found a food source. A common assumption is that Albatrosses must be able to sleep in flight, although no direct evidence has ever been obtained.
This efficient long-distance travelling underlies the albatross's success as a long-distance forager, covering great distances and expending little energy looking for patchily distributed food sources. Their adaptation to gliding flight makes them dependent on wind and waves, but their long wings are ill-suited to powered flight and most species lack the muscles and energy to undertake sustained flapping flight. Albatrosses in calm seas rest on the ocean's surface until the wind picks up again as using powered flight is not energetically worthwhile, though they are capable of flight to avoid danger. The North Pacific albatrosses can use a flight style known as flap-gliding, where the bird progresses by bursts of flapping followed by gliding. When taking off, albatrosses need to take a run up to allow enough air to move under the wing to provide lift.
The dynamic soaring of albatrosses has provided inspiration to airplane designers; German aerospace engineer Johannes Traugott and colleagues have charted the albatross's nuanced flight pattern and are looking for ways to apply this to aircraft, especially in the area of drones and unmanned aircraft.
Distribution and range at sea
Most albatrosses range in the Southern Hemisphere from Antarctica to Australia, South Africa, and South America. The exceptions to this are the four North Pacific albatrosses, of which three occur exclusively in the North Pacific, from Hawaii to Japan, California, and Alaska; and one, the waved albatross, breeds in the Galápagos Islands and feeds off the coast of South America. The need for wind to enable gliding is the reason albatrosses are for the most part confined to higher latitudes; being unsuited to sustained flapping flight makes crossing the doldrums extremely difficult. The exception, the waved albatross, is able to live in the equatorial waters around the Galápagos Islands because of the cool waters of the Humboldt Current and the resulting winds.
Why the albatrosses became extinct in the North Atlantic is unknown for certain, although rising sea levels due to an interglacial warming period are thought to have submerged the site of a short-tailed albatross colony that has been excavated in Bermuda. Some southern species have occasionally turned up as vagrants in the North Atlantic and can become exiled, remaining there for decades. One of these exiles, a black-browed albatross named Albert has been observed travelling to gannet colonies in Scotland for at least 50 years in an attempt to breed. Another black-browed albatross nicknamed Albie has been frequently observed across Northern Europe since 2014, and is also believed to be searching for a mate, having been recorded from Germany, Scandinavia and RSPB Bempton Cliffs in Yorkshire, England.
The use of satellite tracking is teaching scientists a great deal about the way albatrosses range across the ocean to find food. They undertake no annual migration, but disperse widely after breeding; Southern Hemisphere species often undertake circumpolar trips. Evidence also exists of separate ranges for different species at sea. A comparison of the foraging niches of two related species that breed on Campbell Island, the Campbell albatross and the grey-headed albatross, showed the Campbell albatross primarily fed over the Campbell Plateau, whereas the grey-headed albatross fed in more pelagic, oceanic waters. Wandering albatrosses also react strongly to bathymetry, feeding only in waters deeper than ; so rigidly did the satellite plots match this contour that one scientist remarked, "It almost appears as if the birds notice and obey a 'No Entry' sign where the water shallows to less than 1000 (metres)". Also, evidence shows different ranges for the two sexes of the same species; a study of Tristan albatrosses breeding on Gough Island showed that males foraged to the west of Gough and females to the east.
Ecology
Diet
The albatross diet is predominantly cephalopods, fish, crustaceans, and offal (organ meat), although they also scavenge carrion and feed on other zooplankton. For most species, a comprehensive understanding of diet is known for only the breeding season, when the albatrosses regularly return to land and study is possible. The importance of each of these food sources varies from species to species, and even from population to population; some concentrate on squid alone, others take more krill or fish. Of the two albatross species found in Hawaii, one, the black-footed albatross, takes mostly fish, while the Laysan feeds on squid.
The use of data loggers at sea that record ingestion of water against time (providing a likely time of feeding) suggests that albatrosses predominantly feed during the day. Analysis of the squid beaks regurgitated by albatrosses has shown that many of the squid eaten are too large to have been caught alive, and include midwater species likely to be beyond the reach of albatross, suggesting that, for some species (like the wandering albatross), scavenged squid may be an important part of the diet. The source of these dead squid is a matter of debate; some certainly comes from squid fisheries, but in nature it primarily comes from the die-off that occurs after squid spawning and the vomit of squid-eating whales (sperm whales, pilot whales, and southern bottlenose whales). The diet of other species, like the black-browed albatross or the grey-headed albatross, is rich with smaller species of squid that tend to sink after death, and scavenging is not assumed to play a large role in their diet. The waved albatross has been observed practising kleptoparasitism, harassing boobies to steal their food, making it the only member of its order to do so regularly.
Until recently, albatrosses were thought to be predominantly surface feeders, swimming at the surface and snapping up squid and fish pushed to the surface by currents, predators, or death. The deployment of capillary depth recorders, which record the maximum dive depth undertaken by a bird, has shown that while some species, such as the wandering albatross, do not dive deeper than a metre, some species, such as the light-mantled albatross, have a mean diving depth of almost and can dive as deep as . In addition to surface feeding and diving, they have also been observed plunge diving from the air to snatch prey.
Breeding and dancing
Albatrosses are colonial, usually nesting on isolated islands; where colonies are on larger landmasses, they are found on exposed headlands with good approaches from the sea in several directions, like the colony on the Otago Peninsula in Dunedin, New Zealand. Many Buller's albatrosses and black-footed albatrosses nest under trees in open forest. Colonies vary from the very dense aggregations favoured by the mollymawks (black-browed albatross colonies on the Falkland Islands have densities of 70 nests per 100 m2) to the much looser groups and widely spaced individual nests favoured by the sooty and great albatrosses. All albatross colonies are on islands that historically were free of land mammals. Albatrosses are highly philopatric, meaning they usually return to their natal colony to breed. This tendency is so strong that a study of Laysan albatrosses showed that the average distance between hatching site and the site where a bird established its own territory was .
Albatrosses live much longer than other birds; they delay breeding for longer and invest more effort into fewer young. Most species survive upwards of 50 years, the oldest recorded being a Laysan albatross named Wisdom that was ringed in 1956 as a mature adult and hatched another chick in February 2021, making her at least 70 years old. She is the oldest confirmed wild bird and the oldest banded bird in the world.
Albatrosses reach sexual maturity slowly, after about five years, but even once they have reached maturity, they do not begin to breed for another few years (even up to 10 years for some species). Young nonbreeders attend a colony prior to beginning to breed, spending many years practising the elaborate breeding rituals and "dances" for which the family is famous. Birds arriving back at the colony for the first time already have the stereotyped behaviours that compose albatross language, but can neither "read" that behaviour as exhibited by other birds nor respond appropriately.
The repertoire of behaviour involves synchronised performances of various actions such as preening, pointing, calling, bill clacking, staring, and combinations of such behaviours (such as the sky-call).
Albatrosses are held to undertake these elaborate and painstaking rituals to ensure that the appropriate partner has been chosen and to perfect partner recognition, as egg laying and chick rearing is a huge investment. Even species that can complete an egg-laying cycle in under a year seldom lay eggs in consecutive years. The great albatrosses (i.e., wandering albatross) take over a year to raise a chick from laying to fledging. Albatrosses lay a single subelliptical egg, white with reddish-brown spots, in a breeding season; if the egg is lost to predators or accidentally broken, then no further breeding attempts are made that year. The larger eggs weigh from . The "divorce" of a pair is a rare occurrence, due to the diminished lifetime reproductive success it causes, and usually happens only after several years of breeding failure.
All the southern albatrosses create large nests for their egg, using grass, shrubs, soil, peat, and even penguin feathers, whereas the three species in the North Pacific make more rudimentary nests. The waved albatross, though, makes no nest and even moves its egg around the pair's territory, as much as , sometimes causing it to lose the egg. In all albatross species, both parents incubate the egg in stints that last between one day and three weeks. Incubation lasts around 70 to 80 days (longer for the larger albatrosses), the longest incubation period of any bird. It can be an energetically demanding process, with the adult losing as much as of body weight a day.
After hatching, the chick, which is semi-altricial, is brooded and guarded for three weeks until it is large enough to defend and thermoregulate itself. During this period, the parents feed the chick small meals when they relieve each other from duty. After the brooding period is over, the chick is fed in regular intervals by both parents. The parents adopt alternative patterns of short and long foraging trips, providing meals that weigh around 12% of their body weight (around 600 g, or 21 oz). The meals are composed of fresh squid, fish, and krill, as well as stomach oil, an energy-rich food that is lighter to carry than undigested prey items. This oil is created in a stomach organ known as a proventriculus from digested prey items by most Procellariiformes, and gives them their distinctive musty smell.
Albatross chicks take a long time to fledge. In the case of the great albatrosses, it can take up to 280 days; even for the smaller albatrosses, it takes between 140 and 170 days. Like many seabirds, albatross chicks will gain enough weight to be heavier than their parents, and prior to fledging, they use these reserves to build up body condition (particularly growing all their flight feathers), usually fledging at the same weight as their parents. Between 15 and 65% of those fledged survive to breed. Albatross chicks fledge on their own and receive no further help from their parents, which return to the nest after fledging, unaware their chick has left. Studies of juveniles dispersing at sea have suggested an innate migration behaviour, a genetically coded navigation route, which helps young birds when they are first out at sea.
Hybridization is rare in albatrosses, largely due to the low incidence of breeding-site vagrancy.
In culture
Albatrosses are described in the Handbook of World Birds as "the most legendary of all birds". An albatross is the central emblem in The Rime of the Ancient Mariner by Samuel Taylor Coleridge, representing the innocence and beauty of God's creation. The albatross metaphor is derived from this poem; someone bearing a burden or facing an obstacle is said to have "an albatross around his neck", the punishment given to the mariner who killed the albatross. A widespread myth holds that sailors believe shooting or harming an albatross is disastrous, due in part to the poem; in truth, sailors regularly killed and ate them, as reported by James Cook in 1772. However, other sailors reportedly caught the birds but let them free again, possibly believing that albatrosses were the souls of lost sailors, so killing them would bring bad luck. The head of an albatross being caught with a hook is used as the emblem of the Cape Horners, i.e., sailors who have rounded Cape Horn on freighters under sail; captains of such ships even received themselves the title "albatrosses" in the Cape Horners' organisation.
A captive albatross tormented by jeering sailors is also a metaphor for the social travails of the sensitive poète maudit in Charles Baudelaire's poem L'albatros:
In golf, shooting three under par on a single hole has been termed scoring an "albatross", as a continuation on the birdie and eagle theme.
Non-European mythologies
The Māori used the wing bones of the albatross to carve flutes. In Hawaiian mythology, Laysan albatrosses are considered aumakua, being a sacred manifestation of the ancestors, and quite possibly also the sacred bird of Kāne. Japanese mythology, by contrast, refers to the short-tailed albatross as ahodori, "fool bird", due to its habit of disregarding terrestrial predators, making it easy prey for feather collectors.
Birdwatching
Albatrosses are popular birds for birdwatchers, and their colonies are popular destinations for ecotourists. Regular birdwatching trips are taken out of many coastal towns and cities, such as Monterey, Dunedin, Kaikōura, Wollongong, Sydney, Port Fairy, Hobart, and Cape Town, to see pelagic seabirds. Albatrosses are easily attracted to these sightseeing boats by the deployment of fish oil and burley into the sea. Visits to colonies can be very popular; the northern royal albatross colony at Taiaroa Head in Dunedin, New Zealand, attracts 40,000 visitors a year.
Threats and conservation
In spite of often being accorded legendary status, albatrosses have not escaped either indirect or direct pressure from humans. Early encounters with albatrosses by Polynesians and Aleuts resulted in hunting and in some cases extirpation from some islands (such as Easter Island). As Europeans began sailing the world, they, too, began to hunt albatross, "fishing" for them from boats to serve at the table or blasting them for sport. This sport reached its peak on emigration lines bound for Australia, and only died down when ships became too fast to fish from, and regulations forbade the discharge of weapons for safety reasons. In the 19th century, albatross colonies, particularly those in the North Pacific, were harvested for the feather trade, leading to the near-extinction of the short-tailed albatross.
Of the 22 albatross species recognised by IUCN on their Red List, 15 are threatened with extinction, that is, Critically Endangered (Tristan albatross and waved albatross), Endangered (7 species), or Vulnerable (6 species). Six further species are considered as Near Threatened and only one of Least Concern. One of the main threats is commercial longline fishing, as the albatrosses and other seabirds—which will readily feed on offal—are attracted to the set bait, become hooked on the lines and drown. An estimated 100,000 albatross per year are killed in this fashion. Unregulated pirate fisheries exacerbate the problem. A study showed that potentially illegal longline fishing activities are highly concentrated in areas of illegally-caught fish species, and the risk to bycatch albatrosses is significantly higher in areas where these illegal longline fishing vessels operate.
On Midway Atoll, collisions between Laysan albatrosses and aircraft have resulted in human and bird deaths, as well as severe disruptions in military flight operations. Studies were made in the late 1950s and early 1960s that examined the results of control methods such as the killing of birds, the levelling and clearing of land to eliminate updrafts, and the destruction of annual nesting sites. Tall structures such as traffic control and radio towers killed 3000 birds in flight collisions during 1964–1965 before the towers were taken down. Closure of Naval Air Facility Midway in 1993 eliminated the problem of collisions with military aircraft. By 2008, it was noted that reduction in human activity had helped reduce bird deaths, though lead paint pollution near military buildings continued to poison birds by ingestion. Albatross plumes were popular in the early 20th century. In 1909 alone, over 300,000 albatrosses were killed on Midway Island and Laysan Island for their plumes.
Another threat to albatrosses is introduced species, such as rats or feral cats, which directly attack albatrosses or their chicks and eggs. Albatrosses have evolved to breed on islands where land mammals are absent and have not developed defences against them. Even species as small as mice can be detrimental; on Gough Island, the chicks of Tristan albatrosses are attacked and eaten alive by introduced house mice. Introduced species can have other indirect effects; cattle overgrazed essential cover on Amsterdam Island, threatening the Amsterdam albatross; on other islands, introduced plants reduce potential nesting habitat.
Ingestion of plastic flotsam is another problem, one faced by many seabirds. The amount of plastic in the seas has increased dramatically since the first record in the 1960s, coming from waste discarded by ships, offshore dumping, litter on beaches, and waste washed to sea by rivers. It is impossible to digest and takes up space in the stomach or gizzard that should be used for food, or can cause an obstruction that starves the bird directly. Studies of birds in the North Pacific have shown that ingestion of plastics results in declining body weight and body condition. This plastic is sometimes regurgitated and fed to chicks; a study of Laysan albatross chicks on Midway Atoll showed large amounts of ingested plastic in naturally dead chicks compared to healthy chicks killed in accidents. While not the direct cause of death, this plastic causes physiological stress and causes the chick to feel full during feedings, reducing its food intake and the chances of survival.
Scientists and conservationists (most importantly BirdLife International and their partners, who run the Save the Albatross campaign) are working with governments and fishermen to find solutions to the threats albatrosses face. Techniques such as setting longline bait at night, dyeing the bait blue, setting the bait underwater, increasing the amount of weight on lines, and using bird scarers can all reduce the seabird bycatch. For example, a collaborative study between scientists and fishermen in New Zealand successfully tested an underwater setting device for longliners, which set the lines below the reach of vulnerable albatross species. The use of some of these techniques in the Patagonian toothfish fishery in the Falkland Islands is thought to have reduced the number of black-browed albatrosses taken by the fleet between 1994 and 2004. Conservationists have also worked on the field of island restoration, removing introduced species that threaten native wildlife, which protects albatrosses from introduced predators.
One important step towards protecting albatrosses and other seabirds is the 2001 treaty, the Agreement on the Conservation of Albatrosses and Petrels, which came into force in 2004 and has been ratified by thirteen countries, Argentina, Australia, Brazil, Chile, Ecuador, France, New Zealand, Norway, Peru, South Africa, Spain, the United Kingdom, and Uruguay. The treaty requires these countries to take specific actions to reduce bycatch, pollution and to remove introduced species from nesting islands.
Species
Since 1996, albatrosses have been divided into four genera. The number of species is a matter of debate. The IUCN and BirdLife International recognise 22 extant species (listed below), ITIS recognise 21 (the 22 below minus T. steadi), and a 2004 paper proposed a reduction to 13 (indicated in parentheses below), comprising the traditional 14 species minus D. amsterdamensis.
| Biology and health sciences | Procellariiformes | null |
27051151 | https://en.wikipedia.org/wiki/Big%20data | Big data | Big data primarily refers to data sets that are too large or complex to be dealt with by traditional data-processing software. Data with many entries (rows) offer greater statistical power, while data with higher complexity (more attributes or columns) may lead to a higher false discovery rate.
Big data analysis challenges include capturing data, data storage, data analysis, search, sharing, transfer, visualization, querying, updating, information privacy, and data source. Big data was originally associated with three key concepts: volume, variety, and velocity. The analysis of big data presents challenges in sampling, and thus previously allowing for only observations and sampling. Thus a fourth concept, veracity, refers to the quality or insightfulness of the data. Without sufficient investment in expertise for big data veracity, the volume and variety of data can produce costs and risks that exceed an organization's capacity to create and capture value from big data.
Current usage of the term big data tends to refer to the use of predictive analytics, user behavior analytics, or certain other advanced data analytics methods that extract value from big data, and seldom to a particular size of data set. "There is little doubt that the quantities of data now available are indeed large, but that's not the most relevant characteristic of this new data ecosystem."
Analysis of data sets can find new correlations to "spot business trends, prevent diseases, combat crime and so on". Scientists, business executives, medical practitioners, advertising and governments alike regularly meet difficulties with large data-sets in areas including Internet searches, fintech, healthcare analytics, geographic information systems, urban informatics, and business informatics. Scientists encounter limitations in e-Science work, including meteorology, genomics, connectomics, complex physics simulations, biology, and environmental research.
The size and number of available data sets have grown rapidly as data is collected by devices such as mobile devices, cheap and numerous information-sensing Internet of things devices, aerial (remote sensing) equipment, software logs, cameras, microphones, radio-frequency identification (RFID) readers and wireless sensor networks. The world's technological per-capita capacity to store information has roughly doubled every 40 months since the 1980s; , every day 2.5 exabytes (2.17×260 bytes) of data are generated. Based on an IDC report prediction, the global data volume was predicted to grow exponentially from 4.4 zettabytes to 44 zettabytes between 2013 and 2020. By 2025, IDC predicts there will be 163 zettabytes of data. According to IDC, global spending on big data and business analytics (BDA) solutions is estimated to reach $215.7 billion in 2021. While Statista report, the global big data market is forecasted to grow to $103 billion by 2027. In 2011 McKinsey & Company reported, if US healthcare were to use big data creatively and effectively to drive efficiency and quality, the sector could create more than $300 billion in value every year. In the developed economies of Europe, government administrators could save more than €100 billion ($149 billion) in operational efficiency improvements alone by using big data. And users of services enabled by personal-location data could capture $600 billion in consumer surplus. One question for large enterprises is determining who should own big-data initiatives that affect the entire organization.
Relational database management systems and desktop statistical software packages used to visualize data often have difficulty processing and analyzing big data. The processing and analysis of big data may require "massively parallel software running on tens, hundreds, or even thousands of servers". What qualifies as "big data" varies depending on the capabilities of those analyzing it and their tools. Furthermore, expanding capabilities make big data a moving target. "For some organizations, facing hundreds of gigabytes of data for the first time may trigger a need to reconsider data management options. For others, it may take tens or hundreds of terabytes before data size becomes a significant consideration."
Definition
The term big data has been in use since the 1990s, with some giving credit to John Mashey for popularizing the term. Big data usually includes data sets with sizes beyond the ability of commonly used software tools to capture, curate, manage, and process data within a tolerable elapsed time. Big data philosophy encompasses unstructured, semi-structured and structured data; however, the main focus is on unstructured data. Big data "size" is a constantly moving target; ranging from a few dozen terabytes to many zettabytes of data. Big data requires a set of techniques and technologies with new forms of integration to reveal insights from data-sets that are diverse, complex, and of a massive scale.
"Volume", "variety", "velocity", and various other "Vs" are added by some organizations to describe it, a revision challenged by some industry authorities. The Vs of big data were often referred to as the "three Vs", "four Vs", and "five Vs". They represented the qualities of big data in volume, variety, velocity, veracity, and value. Variability is often included as an additional quality of big data.
A 2018 definition states "Big data is where parallel computing tools are needed to handle data", and notes, "This represents a distinct and clearly defined change in the computer science used, via parallel programming theories, and losses of some of the guarantees and capabilities made by Codd's relational model."
In a comparative study of big datasets, Kitchin and McArdle found that none of the commonly considered characteristics of big data appear consistently across all of the analyzed cases. For this reason, other studies identified the redefinition of power dynamics in knowledge discovery as the defining trait. Instead of focusing on the intrinsic characteristics of big data, this alternative perspective pushes forward a relational understanding of the object claiming that what matters is the way in which data is collected, stored, made available and analyzed.
Big data vs. business intelligence
The growing maturity of the concept more starkly delineates the difference between "big data" and "business intelligence":
Business intelligence uses applied mathematics tools and descriptive statistics with data with high information density to measure things, detect trends, etc.
Big data uses mathematical analysis, optimization, inductive statistics, and concepts from nonlinear system identification to infer laws (regressions, nonlinear relationships, and causal effects) from large sets of data with low information density to reveal relationships and dependencies, or to perform predictions of outcomes and behaviors.
Characteristics
Big data can be described by the following characteristics:
Volume The quantity of generated and stored data. The size of the data determines the value and potential insight, and whether it can be considered big data or not. The size of big data is usually larger than terabytes and petabytes.
Variety The type and nature of the data. Earlier technologies like RDBMSs were capable to handle structured data efficiently and effectively. However, the change in type and nature from structured to semi-structured or unstructured challenged the existing tools and technologies. Big data technologies evolved with the prime intention to capture, store, and process the semi-structured and unstructured (variety) data generated with high speed (velocity), and huge in size (volume). Later, these tools and technologies were explored and used for handling structured data also but preferable for storage. Eventually, the processing of structured data was still kept as optional, either using big data or traditional RDBMSs. This helps in analyzing data towards effective usage of the hidden insights exposed from the data collected via social media, log files, sensors, etc. Big data draws from text, images, audio, video; plus it completes missing pieces through data fusion.
Velocity The speed at which the data is generated and processed to meet the demands and challenges that lie in the path of growth and development. Big data is often available in real-time. Compared to small data, big data is produced more continually. Two kinds of velocity related to big data are the frequency of generation and the frequency of handling, recording, and publishing.
Veracity The truthfulness or reliability of the data, which refers to the data quality and the data value. Big data must not only be large in size, but also must be reliable in order to achieve value in the analysis of it. The data quality of captured data can vary greatly, affecting an accurate analysis.
Value The worth in information that can be achieved by the processing and analysis of large datasets. Value also can be measured by an assessment of the other qualities of big data. Value may also represent the profitability of information that is retrieved from the analysis of big data.
Variability The characteristic of the changing formats, structure, or sources of big data. Big data can include structured, unstructured, or combinations of structured and unstructured data. Big data analysis may integrate raw data from multiple sources. The processing of raw data may also involve transformations of unstructured data to structured data.
Other possible characteristics of big data are:
Exhaustive Whether the entire system (i.e., =all) is captured or recorded or not. Big data may or may not include all the available data from sources.
Fine-grained and uniquely lexical Respectively, the proportion of specific data of each element per element collected and if the element and its characteristics are properly indexed or identified.
Relational If the data collected contains common fields that would enable a conjoining, or meta-analysis, of different data sets.
Extensional If new fields in each element of the data collected can be added or changed easily.
Scalability If the size of the big data storage system can expand rapidly.
Architecture
Big data repositories have existed in many forms, often built by corporations with a special need. Commercial vendors historically offered parallel database management systems for big data beginning in the 1990s. For many years, WinterCorp published the largest database report.
Teradata Corporation in 1984 marketed the parallel processing DBC 1012 system. Teradata systems were the first to store and analyze 1 terabyte of data in 1992. Hard disk drives were 2.5 GB in 1991 so the definition of big data continuously evolves. Teradata installed the first petabyte class RDBMS based system in 2007. , there are a few dozen petabyte class Teradata relational databases installed, the largest of which exceeds 50 PB. Systems up until 2008 were 100% structured relational data. Since then, Teradata has added semi structured data types including XML, JSON, and Avro.
In 2000, Seisint Inc. (now LexisNexis Risk Solutions) developed a C++-based distributed platform for data processing and querying known as the HPCC Systems platform. This system automatically partitions, distributes, stores and delivers structured, semi-structured, and unstructured data across multiple commodity servers. Users can write data processing pipelines and queries in a declarative dataflow programming language called ECL. Data analysts working in ECL are not required to define data schemas upfront and can rather focus on the particular problem at hand, reshaping data in the best possible manner as they develop the solution. In 2004, LexisNexis acquired Seisint Inc. and their high-speed parallel processing platform and successfully used this platform to integrate the data systems of Choicepoint Inc. when they acquired that company in 2008. In 2011, the HPCC systems platform was open-sourced under the Apache v2.0 License.
CERN and other physics experiments have collected big data sets for many decades, usually analyzed via high-throughput computing rather than the map-reduce architectures usually meant by the current "big data" movement.
In 2004, Google published a paper on a process called MapReduce that uses a similar architecture. The MapReduce concept provides a parallel processing model, and an associated implementation was released to process huge amounts of data. With MapReduce, queries are split and distributed across parallel nodes and processed in parallel (the "map" step). The results are then gathered and delivered (the "reduce" step). The framework was very successful, so others wanted to replicate the algorithm. Therefore, an implementation of the MapReduce framework was adopted by an Apache open-source project named "Hadoop". Apache Spark was developed in 2012 in response to limitations in the MapReduce paradigm, as it adds in-memory processing and the ability to set up many operations (not just map followed by reducing).
MIKE2.0 is an open approach to information management that acknowledges the need for revisions due to big data implications identified in an article titled "Big Data Solution Offering". The methodology addresses handling big data in terms of useful permutations of data sources, complexity in interrelationships, and difficulty in deleting (or modifying) individual records.
Studies in 2012 showed that a multiple-layer architecture was one option to address the issues that big data presents. A distributed parallel architecture distributes data across multiple servers; these parallel execution environments can dramatically improve data processing speeds. This type of architecture inserts data into a parallel DBMS, which implements the use of MapReduce and Hadoop frameworks. This type of framework looks to make the processing power transparent to the end-user by using a front-end application server.
The data lake allows an organization to shift its focus from centralized control to a shared model to respond to the changing dynamics of information management. This enables quick segregation of data into the data lake, thereby reducing the overhead time.
Technologies
A 2011 McKinsey Global Institute report characterizes the main components and ecosystem of big data as follows:
Techniques for analyzing data, such as A/B testing, machine learning, and natural language processing
Big data technologies, like business intelligence, cloud computing, and databases
Visualization, such as charts, graphs, and other displays of the data
Multidimensional big data can also be represented as OLAP data cubes or, mathematically, tensors. Array database systems have set out to provide storage and high-level query support on this data type.
Additional technologies being applied to big data include efficient tensor-based computation, such as multilinear subspace learning, massively parallel-processing (MPP) databases, search-based applications, data mining, distributed file systems, distributed cache (e.g., burst buffer and Memcached), distributed databases, cloud and HPC-based infrastructure (applications, storage and computing resources), and the Internet. Although, many approaches and technologies have been developed, it still remains difficult to carry out machine learning with big data.
Some MPP relational databases have the ability to store and manage petabytes of data. Implicit is the ability to load, monitor, back up, and optimize the use of the large data tables in the RDBMS.
DARPA's Topological Data Analysis program seeks the fundamental structure of massive data sets and in 2008 the technology went public with the launch of a company called "Ayasdi".
The practitioners of big data analytics processes are generally hostile to slower shared storage, preferring direct-attached storage (DAS) in its various forms from solid state drive (SSD) to high capacity SATA disk buried inside parallel processing nodes. The perception of shared storage architectures—storage area network (SAN) and network-attached storage (NAS)— is that they are relatively slow, complex, and expensive. These qualities are not consistent with big data analytics systems that thrive on system performance, commodity infrastructure, and low cost.
Real or near-real-time information delivery is one of the defining characteristics of big data analytics. Latency is therefore avoided whenever and wherever possible. Data in direct-attached memory or disk is good—data on memory or disk at the other end of an FC SAN connection is not. The cost of an SAN at the scale needed for analytics applications is much higher than other storage techniques.
Applications
Big data has increased the demand of information management specialists so much so that Software AG, Oracle Corporation, IBM, Microsoft, SAP, EMC, HP, and Dell have spent more than $15 billion on software firms specializing in data management and analytics. In 2010, this industry was worth more than $100 billion and was growing at almost 10 percent a year, about twice as fast as the software business as a whole.
Developed economies increasingly use data-intensive technologies. There are 4.6 billion mobile-phone subscriptions worldwide, and between 1 billion and 2 billion people accessing the internet. Between 1990 and 2005, more than 1 billion people worldwide entered the middle class, which means more people became more literate, which in turn led to information growth. The world's effective capacity to exchange information through telecommunication networks was 281 petabytes in 1986, 471 petabytes in 1993, 2.2 exabytes in 2000, 65 exabytes in 2007 and predictions put the amount of internet traffic at 667 exabytes annually by 2014. According to one estimate, one-third of the globally stored information is in the form of alphanumeric text and still image data, which is the format most useful for most big data applications. This also shows the potential of yet unused data (i.e. in the form of video and audio content).
While many vendors offer off-the-shelf products for big data, experts promote the development of in-house custom-tailored systems if the company has sufficient technical capabilities.
Government
The use and adoption of big data within governmental processes allows efficiencies in terms of cost, productivity, and innovation, but comes with flaws. Data analysis often requires multiple parts of government (central and local) to work in collaboration and create new and innovative processes to deliver the desired outcome. A common government organization that makes use of big data is the National Security Administration (NSA), which monitors the activities of the Internet constantly in search for potential patterns of suspicious or illegal activities their system may pick up.
Civil registration and vital statistics (CRVS) collects all certificates status from birth to death. CRVS is a source of big data for governments.
International development
Research on the effective usage of information and communication technologies for development (also known as "ICT4D") suggests that big data technology can make important contributions but also present unique challenges to international development. Advancements in big data analysis offer cost-effective opportunities to improve decision-making in critical development areas such as health care, employment, economic productivity, crime, security, and natural disaster and resource management. Additionally, user-generated data offers new opportunities to give the unheard a voice. However, longstanding challenges for developing regions such as inadequate technological infrastructure and economic and human resource scarcity exacerbate existing concerns with big data such as privacy, imperfect methodology, and interoperability issues. The challenge of "big data for development" is currently evolving toward the application of this data through machine learning, known as "artificial intelligence for development (AI4D).
Benefits
A major practical application of big data for development has been "fighting poverty with data". In 2015, Blumenstock and colleagues estimated predicted poverty and wealth from mobile phone metadata and in 2016 Jean and colleagues combined satellite imagery and machine learning to predict poverty. Using digital trace data to study the labor market and the digital economy in Latin America, Hilbert and colleagues argue that digital trace data has several benefits such as:
Thematic coverage: including areas that were previously difficult or impossible to measure
Geographical coverage: providing sizable and comparable data for almost all countries, including many small countries that usually are not included in international inventories
Level of detail: providing fine-grained data with many interrelated variables, and new aspects, like network connections
Timeliness and timeseries: graphs can be produced within days of being collected
Challenges
At the same time, working with digital trace data instead of traditional survey data does not eliminate the traditional challenges involved when working in the field of international quantitative analysis. Priorities change, but the basic discussions remain the same. Among the main challenges are:
Representativeness. While traditional development statistics is mainly concerned with the representativeness of random survey samples, digital trace data is never a random sample.
Generalizability. While observational data always represents this source very well, it only represents what it represents, and nothing more. While it is tempting to generalize from specific observations of one platform to broader settings, this is often very deceptive.
Harmonization. Digital trace data still requires international harmonization of indicators. It adds the challenge of so-called "data-fusion", the harmonization of different sources.
Data overload. Analysts and institutions are not used to effectively deal with a large number of variables, which is efficiently done with interactive dashboards. Practitioners still lack a standard workflow that would allow researchers, users and policymakers to efficiently and effectively deal with data.
Finance
Big Data is being rapidly adopted in Finance to 1) speed up processing and 2) deliver better, more informed inferences, both internally and to the clients of the financial institutions. The financial applications of Big Data range from investing decisions and trading (processing volumes of available price data, limit order books, economic data and more, all at the same time), portfolio management (optimizing over an increasingly large array of financial instruments, potentially selected from different asset classes), risk management (credit rating based on extended information), and any other aspect where the data inputs are large. Big Data has also been a typical concept within the field of alternative financial service. Some of the major areas involve crowd-funding platforms and crypto currency exchanges.
Healthcare
Big data analytics has been used in healthcare in providing personalized medicine and prescriptive analytics, clinical risk intervention and predictive analytics, waste and care variability reduction, automated external and internal reporting of patient data, standardized medical terms and patient registries. Some areas of improvement are more aspirational than actually implemented. The level of data generated within healthcare systems is not trivial. With the added adoption of mHealth, eHealth and wearable technologies the volume of data will continue to increase. This includes electronic health record data, imaging data, patient generated data, sensor data, and other forms of difficult to process data. There is now an even greater need for such environments to pay greater attention to data and information quality. "Big data very often means 'dirty data' and the fraction of data inaccuracies increases with data volume growth." Human inspection at the big data scale is impossible and there is a desperate need in health service for intelligent tools for accuracy and believability control and handling of information missed. While extensive information in healthcare is now electronic, it fits under the big data umbrella as most is unstructured and difficult to use. The use of big data in healthcare has raised significant ethical challenges ranging from risks for individual rights, privacy and autonomy, to transparency and trust.
Big data in health research is particularly promising in terms of exploratory biomedical research, as data-driven analysis can move forward more quickly than hypothesis-driven research. Then, trends seen in data analysis can be tested in traditional, hypothesis-driven follow up biological research and eventually clinical research.
A related application sub-area, that heavily relies on big data, within the healthcare field is that of computer-aided diagnosis in medicine. For instance, for epilepsy monitoring it is customary to create 5 to 10 GB of data daily. Similarly, a single uncompressed image of breast tomosynthesis averages 450 MB of data. These are just a few of the many examples where computer-aided diagnosis uses big data. For this reason, big data has been recognized as one of the seven key challenges that computer-aided diagnosis systems need to overcome in order to reach the next level of performance.
Education
A McKinsey Global Institute study found a shortage of 1.5 million highly trained data professionals and managers and a number of universities including University of Tennessee and UC Berkeley, have created masters programs to meet this demand. Private boot camps have also developed programs to meet that demand, including paid programs like The Data Incubator or General Assembly. In the specific field of marketing, one of the problems stressed by Wedel and Kannan is that marketing has several sub domains (e.g., advertising, promotions,
product development, branding) that all use different types of data.
Media
To understand how the media uses big data, it is first necessary to provide some context into the mechanism used for media process. It has been suggested by Nick Couldry and Joseph Turow that practitioners in media and advertising approach big data as many actionable points of information about millions of individuals. The industry appears to be moving away from the traditional approach of using specific media environments such as newspapers, magazines, or television shows and instead taps into consumers with technologies that reach targeted people at optimal times in optimal locations. The ultimate aim is to serve or convey, a message or content that is (statistically speaking) in line with the consumer's mindset. For example, publishing environments are increasingly tailoring messages (advertisements) and content (articles) to appeal to consumers that have been exclusively gleaned through various data-mining activities.
Targeting of consumers (for advertising by marketers)
Data capture
Data journalism: publishers and journalists use big data tools to provide unique and innovative insights and infographics.
Channel 4, the British public-service television broadcaster, is a leader in the field of big data and data analysis.
Insurance
Health insurance providers are collecting data on social "determinants of health" such as food and TV consumption, marital status, clothing size, and purchasing habits, from which they make predictions on health costs, in order to spot health issues in their clients. It is controversial whether these predictions are currently being used for pricing.
Internet of things (IoT)
Big data and the IoT work in conjunction. Data extracted from IoT devices provides a mapping of device inter-connectivity. Such mappings have been used by the media industry, companies, and governments to more accurately target their audience and increase media efficiency. The IoT is also increasingly adopted as a means of gathering sensory data, and this sensory data has been used in medical, manufacturing and transportation contexts.
Kevin Ashton, the digital innovation expert who is credited with coining the term, defines the Internet of things in this quote: "If we had computers that knew everything there was to know about things—using data they gathered without any help from us—we would be able to track and count everything, and greatly reduce waste, loss, and cost. We would know when things needed replacing, repairing, or recalling, and whether they were fresh or past their best."
Information technology
Especially since 2015, big data has come to prominence within business operations as a tool to help employees work more efficiently and streamline the collection and distribution of information technology (IT). The use of big data to resolve IT and data collection issues within an enterprise is called IT operations analytics (ITOA). By applying big data principles into the concepts of machine intelligence and deep computing, IT departments can predict potential issues and prevent them. ITOA businesses offer platforms for systems management that bring data silos together and generate insights from the whole of the system rather than from isolated pockets of data.
Survey science
Compared to survey-based data collection, big data has low cost per data point, applies analysis techniques via machine learning and data mining, and includes diverse and new data sources, e.g., registers, social media, apps, and other forms digital data. Since 2018, survey scientists have started to examine how big data and survey science can complement each other to allow researchers and practitioners to improve the production of statistics and its quality. There have been three Big Data Meets Survey Science (BigSurv) conferences in 2018, 2020 (virtual), 2023, and one conference forthcoming in 2025, a special issue in the Social Science Computer Review, a special issue in Journal of the Royal Statistical Society, and a special issue in EP J Data Science, and a book called Big Data Meets Social Sciences edited by Craig Hill and five other Fellows of the American Statistical Association. In 2021, the founding members of BigSurv received the Warren J. Mitofsky Innovators Award from the American Association for Public Opinion Research.
Marketing
Big data is notable in marketing due to the constant "datafication" of everyday consumers of the internet, in which all forms of data are tracked. The datafication of consumers can be defined as quantifying many of or all human behaviors for the purpose of marketing. The increasingly digital world of rapid datafication makes this idea relevant to marketing because the amount of data constantly grows exponentially. It is predicted to increase from 44 to 163 zettabytes within the span of five years. The size of big data can often be difficult to navigate for marketers. As a result, adopters of big data may find themselves at a disadvantage. Algorithmic findings can be difficult to achieve with such large datasets. Big data in marketing is a highly lucrative tool that can be used for large corporations, its value being as a result of the possibility of predicting significant trends, interests, or statistical outcomes in a consumer-based manner.
There are three significant factors in the use of big data in marketing:
Big data provides customer behavior pattern spotting for marketers, since all human actions are being quantified into readable numbers for marketers to analyze and use for their research. In addition, big data can also be seen as a customized product recommendation tool. Specifically, since big data is effective in analyzing customers' purchase behaviors and browsing patterns, this technology can assist companies in promoting specific personalized products to specific customers.
Real-time market responsiveness is important for marketers because of the ability to shift marketing efforts and correct to current trends, which is helpful in maintaining relevance to consumers. This can supply corporations with the information necessary to predict the wants and needs of consumers in advance.
Data-driven market ambidexterity are being highly fueled by big data. New models and algorithms are being developed to make significant predictions about certain economic and social situations.
Case studies
Government
China
The Integrated Joint Operations Platform (IJOP, 一体化联合作战平台) is used by the government to monitor the population, particularly Uyghurs. Biometrics, including DNA samples, are gathered through a program of free physicals.
By 2020, China plans to give all its citizens a personal "social credit" score based on how they behave. The Social Credit System, now being piloted in a number of Chinese cities, is considered a form of mass surveillance which uses big data analysis technology.
India
Big data analysis was tried out for the BJP to win the 2014 Indian General Election.
The Indian government uses numerous techniques to ascertain how the Indian electorate is responding to government action, as well as ideas for policy augmentation.
Israel
Personalized diabetic treatments can be created through GlucoMe's big data solution.
United Kingdom
Examples of uses of big data in public services:
Data on prescription drugs: by connecting origin, location and the time of each prescription, a research unit was able to exemplify and examine the considerable delay between the release of any given drug, and a UK-wide adaptation of the National Institute for Health and Care Excellence guidelines. This suggests that new or most up-to-date drugs take some time to filter through to the general patient.
Joining up data: a local authority blended data about services, such as road gritting rotas, with services for people at risk, such as Meals on Wheels. The connection of data allowed the local authority to avoid any weather-related delay.
United States
In 2012, the Obama administration announced the Big Data Research and Development Initiative, to explore how big data could be used to address important problems faced by the government. The initiative is composed of 84 different big data programs spread across six departments.
Big data analysis played a large role in Barack Obama's successful 2012 re-election campaign.
The United States Federal Government owns four of the ten most powerful supercomputers in the world.
The Utah Data Center has been constructed by the United States National Security Agency. When finished, the facility will be able to handle a large amount of information collected by the NSA over the Internet. The exact amount of storage space is unknown, but more recent sources claim it will be on the order of a few exabytes. This has posed security concerns regarding the anonymity of the data collected.
Retail
Walmart handles more than 1 million customer transactions every hour, which are imported into databases estimated to contain more than 2.5 petabytes (2560 terabytes) of data—the equivalent of 167 times the information contained in all the books in the US Library of Congress.
Windermere Real Estate uses location information from nearly 100 million drivers to help new home buyers determine their typical drive times to and from work throughout various times of the day.
FICO Card Detection System protects accounts worldwide.
Omnichannel retailing leverages online big data to improve offline experiences.
Science
The Large Hadron Collider experiments represent about 150 million sensors delivering data 40 million times per second. There are nearly 600 million collisions per second. After filtering and refraining from recording more than 99.99995% of these streams, there are 1,000 collisions of interest per second.
As a result, only working with less than 0.001% of the sensor stream data, the data flow from all four LHC experiments represents 25 petabytes annual rate before replication (). This becomes nearly 200 petabytes after replication.
If all sensor data were recorded in LHC, the data flow would be extremely hard to work with. The data flow would exceed 150 million petabytes annual rate, or nearly 500 exabytes per day, before replication. To put the number in perspective, this is equivalent to 500 quintillion (5×1020) bytes per day, almost 200 times more than all the other sources combined in the world.
The Square Kilometre Array is a radio telescope built of thousands of antennas. It is expected to be operational by 2024. Collectively, these antennas are expected to gather 14 exabytes and store one petabyte per day. It is considered one of the most ambitious scientific projects ever undertaken.
When the Sloan Digital Sky Survey (SDSS) began to collect astronomical data in 2000, it amassed more in its first few weeks than all data collected in the history of astronomy previously. Continuing at a rate of about 200 GB per night, SDSS has amassed more than 140 terabytes of information. When the Large Synoptic Survey Telescope, successor to SDSS, comes online in 2020, its designers expect it to acquire that amount of data every five days.
Decoding the human genome originally took 10 years to process; now it can be achieved in less than a day. The DNA sequencers have divided the sequencing cost by 10,000 in the last ten years, which is 100 times less expensive than the reduction in cost predicted by Moore's law.
The NASA Center for Climate Simulation (NCCS) stores 32 petabytes of climate observations and simulations on the Discover supercomputing cluster.
Google's DNAStack compiles and organizes DNA samples of genetic data from around the world to identify diseases and other medical defects. These fast and exact calculations eliminate any "friction points", or human errors that could be made by one of the numerous science and biology experts working with the DNA. DNAStack, a part of Google Genomics, allows scientists to use the vast sample of resources from Google's search server to scale social experiments that would usually take years, instantly.
23andme's DNA database contains the genetic information of over 1,000,000 people worldwide. The company explores selling the "anonymous aggregated genetic data" to other researchers and pharmaceutical companies for research purposes if patients give their consent. Ahmad Hariri, professor of psychology and neuroscience at Duke University who has been using 23andMe in his research since 2009 states that the most important aspect of the company's new service is that it makes genetic research accessible and relatively cheap for scientists. A study that identified 15 genome sites linked to depression in 23andMe's database lead to a surge in demands to access the repository with 23andMe fielding nearly 20 requests to access the depression data in the two weeks after publication of the paper.
Computational fluid dynamics (CFD) and hydrodynamic turbulence research generate massive data sets. The Johns Hopkins Turbulence Databases (JHTDB) contains over 350 terabytes of spatiotemporal fields from Direct Numerical simulations of various turbulent flows. Such data have been difficult to share using traditional methods such as downloading flat simulation output files. The data within JHTDB can be accessed using "virtual sensors" with various access modes ranging from direct web-browser queries, access through Matlab, Python, Fortran and C programs executing on clients' platforms, to cut out services to download raw data. The data have been used in over 150 scientific publications.
Sports
Big data can be used to improve training and understanding competitors, using sport sensors. It is also possible to predict winners in a match using big data analytics.
Future performance of players could be predicted as well. Thus, players' value and salary is determined by data collected throughout the season.
In Formula One races, race cars with hundreds of sensors generate terabytes of data. These sensors collect data points from tire pressure to fuel burn efficiency.
Based on the data, engineers and data analysts decide whether adjustments should be made in order to win a race. Besides, using big data, race teams try to predict the time they will finish the race beforehand, based on simulations using data collected over the season.
Technology
, eBay.com uses two data warehouses at 7.5 petabytes and 40PB as well as a 40PB Hadoop cluster for search, consumer recommendations, and merchandising.
Amazon.com handles millions of back-end operations every day, as well as queries from more than half a million third-party sellers. The core technology that keeps Amazon running is Linux-based and they had the world's three largest Linux databases, with capacities of 7.8 TB, 18.5 TB, and 24.7 TB.
Facebook handles 50 billion photos from its user base. , Facebook reached 2 billion monthly active users.
Google was handling roughly 100 billion searches per month .
COVID-19
During the COVID-19 pandemic, big data was raised as a way to minimise the impact of the disease. Significant applications of big data included minimising the spread of the virus, case identification and development of medical treatment.
Governments used big data to track infected people to minimise spread. Early adopters included China, Taiwan, South Korea, and Israel.
Research activities
Encrypted search and cluster formation in big data were demonstrated in March 2014 at the American Society of Engineering Education. Gautam Siwach engaged at Tackling the challenges of Big Data by MIT Computer Science and Artificial Intelligence Laboratory and Amir Esmailpour at the UNH Research Group investigated the key features of big data as the formation of clusters and their interconnections. They focused on the security of big data and the orientation of the term towards the presence of different types of data in an encrypted form at cloud interface by providing the raw definitions and real-time examples within the technology. Moreover, they proposed an approach for identifying the encoding technique to advance towards an expedited search over encrypted text leading to the security enhancements in big data.
In March 2012, The White House announced a national "Big Data Initiative" that consisted of six federal departments and agencies committing more than $200 million to big data research projects.
The initiative included a National Science Foundation "Expeditions in Computing" grant of $10 million over five years to the AMPLab at the University of California, Berkeley. The AMPLab also received funds from DARPA, and over a dozen industrial sponsors and uses big data to attack a wide range of problems from predicting traffic congestion to fighting cancer.
The White House Big Data Initiative also included a commitment by the Department of Energy to provide $25 million in funding over five years to establish the Scalable Data Management, Analysis and Visualization (SDAV) Institute, led by the Energy Department's Lawrence Berkeley National Laboratory. The SDAV Institute aims to bring together the expertise of six national laboratories and seven universities to develop new tools to help scientists manage and visualize data on the department's supercomputers.
The U.S. state of Massachusetts announced the Massachusetts Big Data Initiative in May 2012, which provides funding from the state government and private companies to a variety of research institutions. The Massachusetts Institute of Technology hosts the Intel Science and Technology Center for Big Data in the MIT Computer Science and Artificial Intelligence Laboratory, combining government, corporate, and institutional funding and research efforts.
The European Commission is funding the two-year-long Big Data Public Private Forum through their Seventh Framework Program to engage companies, academics and other stakeholders in discussing big data issues. The project aims to define a strategy in terms of research and innovation to guide supporting actions from the European Commission in the successful implementation of the big data economy. Outcomes of this project will be used as input for Horizon 2020, their next framework program.
The British government announced in March 2014 the founding of the Alan Turing Institute, named after the computer pioneer and code-breaker, which will focus on new ways to collect and analyze large data sets.
At the University of Waterloo Stratford Campus Canadian Open Data Experience (CODE) Inspiration Day, participants demonstrated how using data visualization can increase the understanding and appeal of big data sets and communicate their story to the world.
Computational social sciences – Anyone can use application programming interfaces (APIs) provided by big data holders, such as Google and Twitter, to do research in the social and behavioral sciences. Often these APIs are provided for free. Tobias Preis et al. used Google Trends data to demonstrate that Internet users from countries with a higher per capita gross domestic products (GDPs) are more likely to search for information about the future than information about the past. The findings suggest there may be a link between online behaviors and real-world economic indicators. The authors of the study examined Google queries logs made by ratio of the volume of searches for the coming year (2011) to the volume of searches for the previous year (2009), which they call the "future orientation index". They compared the future orientation index to the per capita GDP of each country, and found a strong tendency for countries where Google users inquire more about the future to have a higher GDP.
Tobias Preis and his colleagues Helen Susannah Moat and H. Eugene Stanley introduced a method to identify online precursors for stock market moves, using trading strategies based on search volume data provided by Google Trends. Their analysis of Google search volume for 98 terms of varying financial relevance, published in Scientific Reports, suggests that increases in search volume for financially relevant search terms tend to precede large losses in financial markets.
Big data sets come with algorithmic challenges that previously did not exist. Hence, there is seen by some to be a need to fundamentally change the processing ways.
Sampling big data
A research question that is asked about big data sets is whether it is necessary to look at the full data to draw certain conclusions about the properties of the data or if is a sample is good enough. The name big data itself contains a term related to size and this is an important characteristic of big data. But sampling enables the selection of right data points from within the larger data set to estimate the characteristics of the whole population. In manufacturing different types of sensory data such as acoustics, vibration, pressure, current, voltage, and controller data are available at short time intervals. To predict downtime it may not be necessary to look at all the data but a sample may be sufficient. Big data can be broken down by various data point categories such as demographic, psychographic, behavioral, and transactional data. With large sets of data points, marketers are able to create and use more customized segments of consumers for more strategic targeting.
Critique
Critiques of the big data paradigm come in two flavors: those that question the implications of the approach itself, and those that question the way it is currently done. One approach to this criticism is the field of critical data studies.
Critiques of the big data paradigm
"A crucial problem is that we do not know much about the underlying empirical micro-processes that lead to the emergence of the[se] typical network characteristics of Big Data." In their critique, Snijders, Matzat, and Reips point out that often very strong assumptions are made about mathematical properties that may not at all reflect what is really going on at the level of micro-processes. Mark Graham has leveled broad critiques at Chris Anderson's assertion that big data will spell the end of theory: focusing in particular on the notion that big data must always be contextualized in their social, economic, and political contexts. Even as companies invest eight- and nine-figure sums to derive insight from information streaming in from suppliers and customers, less than 40% of employees have sufficiently mature processes and skills to do so. To overcome this insight deficit, big data, no matter how comprehensive or well analyzed, must be complemented by "big judgment", according to an article in the Harvard Business Review.
Much in the same line, it has been pointed out that the decisions based on the analysis of big data are inevitably "informed by the world as it was in the past, or, at best, as it currently is". Fed by a large number of data on past experiences, algorithms can predict future development if the future is similar to the past. If the system's dynamics of the future change (if it is not a stationary process), the past can say little about the future. In order to make predictions in changing environments, it would be necessary to have a thorough understanding of the systems dynamic, which requires theory. As a response to this critique Alemany Oliver and Vayre suggest to use "abductive reasoning as a first step in the research process in order to bring context to consumers' digital traces and make new theories emerge". Additionally, it has been suggested to combine big data approaches with computer simulations, such as agent-based models and complex systems. Agent-based models are increasingly getting better in predicting the outcome of social complexities of even unknown future scenarios through computer simulations that are based on a collection of mutually interdependent algorithms. Finally, the use of multivariate methods that probe for the latent structure of the data, such as factor analysis and cluster analysis, have proven useful as analytic approaches that go well beyond the bi-variate approaches (e.g. contingency tables) typically employed with smaller data sets.
In health and biology, conventional scientific approaches are based on experimentation. For these approaches, the limiting factor is the relevant data that can confirm or refute the initial hypothesis. A new postulate is accepted now in biosciences: the information provided by the data in huge volumes (omics) without prior hypothesis is complementary and sometimes necessary to conventional approaches based on experimentation. In the massive approaches it is the formulation of a relevant hypothesis to explain the data that is the limiting factor. The search logic is reversed and the limits of induction ("Glory of Science and Philosophy scandal", C. D. Broad, 1926) are to be considered.
Privacy advocates are concerned about the threat to privacy represented by increasing storage and integration of personally identifiable information; expert panels have released various policy recommendations to conform practice to expectations of privacy. The misuse of big data in several cases by media, companies, and even the government has allowed for abolition of trust in almost every fundamental institution holding up society.
Barocas and Nissenbaum argue that one way of protecting individual users is by being informed about the types of information being collected, with whom it is shared, under what constraints and for what purposes.
Critiques of the "V" model
The "V" model of big data is concerning as it centers around computational scalability and lacks in a loss around the perceptibility and understandability of information. This led to the framework of cognitive big data, which characterizes big data applications according to:
Data completeness: understanding of the non-obvious from data
Data correlation, causation, and predictability: causality as not essential requirement to achieve predictability
Explainability and interpretability: humans desire to understand and accept what they understand, where algorithms do not cope with this
Level of automated decision-making: algorithms that support automated decision making and algorithmic self-learning
Critiques of novelty
Large data sets have been analyzed by computing machines for well over a century, including the US census analytics performed by IBM's punch-card machines which computed statistics including means and variances of populations across the whole continent. In more recent decades, science experiments such as CERN have produced data on similar scales to current commercial "big data". However, science experiments have tended to analyze their data using specialized custom-built high-performance computing (super-computing) clusters and grids, rather than clouds of cheap commodity computers as in the current commercial wave, implying a difference in both culture and technology stack.
Critiques of big data execution
Ulf-Dietrich Reips and Uwe Matzat wrote in 2014 that big data had become a "fad" in scientific research. Researcher Danah Boyd has raised concerns about the use of big data in science neglecting principles such as choosing a representative sample by being too concerned about handling the huge amounts of data. This approach may lead to results that have a bias in one way or another. Integration across heterogeneous data resources—some that might be considered big data and others not—presents formidable logistical as well as analytical challenges, but many researchers argue that such integrations are likely to represent the most promising new frontiers in science.
In the provocative article "Critical Questions for Big Data", the authors title big data a part of mythology: "large data sets offer a higher form of intelligence and knowledge [...], with the aura of truth, objectivity, and accuracy". Users of big data are often "lost in the sheer volume of numbers", and "working with Big Data is still subjective, and what it quantifies does not necessarily have a closer claim on objective truth". Recent developments in BI domain, such as pro-active reporting especially target improvements in the usability of big data, through automated filtering of non-useful data and correlations. Big structures are full of spurious correlations either because of non-causal coincidences (law of truly large numbers), solely nature of big randomness (Ramsey theory), or existence of non-included factors so the hope, of early experimenters to make large databases of numbers "speak for themselves" and revolutionize scientific method, is questioned. Catherine Tucker has pointed to "hype" around big data, writing "By itself, big data is unlikely to be valuable." The article explains: "The many contexts where data is cheap relative to the cost of retaining talent to process it, suggests that processing skills are more important than data itself in creating value for a firm."
Big data analysis is often shallow compared to analysis of smaller data sets. In many big data projects, there is no large data analysis happening, but the challenge is the extract, transform, load part of data pre-processing.
Big data is a buzzword and a "vague term", but at the same time an "obsession" with entrepreneurs, consultants, scientists, and the media. Big data showcases such as Google Flu Trends failed to deliver good predictions in recent years, overstating the flu outbreaks by a factor of two. Similarly, Academy awards and election predictions solely based on Twitter were more often off than on target.
Big data often poses the same challenges as small data; adding more data does not solve problems of bias, but may emphasize other problems. In particular data sources such as Twitter are not representative of the overall population, and results drawn from such sources may then lead to wrong conclusions. Google Translate—which is based on big data statistical analysis of text—does a good job at translating web pages. However, results from specialized domains may be dramatically skewed.
On the other hand, big data may also introduce new problems, such as the multiple comparisons problem: simultaneously testing a large set of hypotheses is likely to produce many false results that mistakenly appear significant.
Ioannidis argued that "most published research findings are false" due to essentially the same effect: when many scientific teams and researchers each perform many experiments (i.e. process a big amount of scientific data; although not with big data technology), the likelihood of a "significant" result being false grows fast – even more so, when only positive results are published.
Furthermore, big data analytics results are only as good as the model on which they are predicated. In an example, big data took part in attempting to predict the results of the 2016 U.S. presidential election with varying degrees of success.
Critiques of big data policing and surveillance
Big data has been used in policing and surveillance by institutions like law enforcement and corporations. Due to the less visible nature of data-based surveillance as compared to traditional methods of policing, objections to big data policing are less likely to arise. According to Sarah Brayne's Big Data Surveillance: The Case of Policing, big data policing can reproduce existing societal inequalities in three ways:
Placing people under increased surveillance by using the justification of a mathematical and therefore unbiased algorithm
Increasing the scope and number of people that are subject to law enforcement tracking and exacerbating existing racial overrepresentation in the criminal justice system
Encouraging members of society to abandon interactions with institutions that would create a digital trace, thus creating obstacles to social inclusion
If these potential problems are not corrected or regulated, the effects of big data policing may continue to shape societal hierarchies. Conscientious usage of big data policing could prevent individual level biases from becoming institutional biases, Brayne also notes.
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27051968 | https://en.wikipedia.org/wiki/Indigenous%20horticulture | Indigenous horticulture | Indigenous horticulture is practised in various ways across all inhabited continents. Indigenous refers to the native peoples of a given area and horticulture is the practice of small-scale intercropping.
Africa
North Africa
In North Africa, one such example is the farming practices of the Eggon, a Nigerian hill farming community. The Eggon live in the Mada hills, between Lafia and Akwanga. The hills lay between two rivers, the Mada and Arikya. The altitude helps crops retain moisture on the hills, due to early morning mists and fogs; this also makes for earlier and longer crop cultivation. They practice bush fallow agriculture as well as mixed farming land management styles. They focus on growing yams, cassava, maize, beans, and African rice; much of what is produced is exported as a cash crop and is their secondary source of cash income. The Eggon use a terraced agricultural system to maximize space on the hills. The goats they raise are kept mostly for fertilizers used in farming. They are only killed on special occasions, such as weddings. The Eggon use the diversity in their environment to maximize their crop production.
West Africa
In West Africa, the Kissidougou live on the savannah, dotted by dense areas of "forest islands" created by them. The Kissidougou practice intercropping within the forested areas. However, they also operate farms maintained on slopes or plateaus located between the forest islands. They prepare the savannah lands for forestation through farming and burning of the grasses to fertilize the soils. The Kissidougou graze their cattle on the savannah to help to maintain flammable grasses around the farms and the villages. The Kissidougou create diversity in their environment by farming and transforming savannah into lush, dense forest.
The prevalence of wetlands in West Africa has helped to support local indigenous horticulture. Seasonal flooding of major rivers in the region, such as the Niger, the Sudd, and the Senegal, have made it possible for flood-cropping in many areas. Indigenous people have utilized a variety of irrigation techniques in order to take advantage of this flooding. Additionally, they plan their plantings and harvests specifically around the flooding of local rivers. For example, some farmers choose to plant on rising floods and harvest as the flooding diminishes. This techniques is utilized when cultivating rice.
East Africa
East Africa is one of the areas most affected by corporate farming. In western Uganda, there is a farming society called the Banyankore. Their land is part of the Ryampara hill country, between two flat, dry, and less populated areas of land. On the hills, the average monthly rainfall is 970mm, with two short dry months in June and July and a short, less severe, dry month in January. Seven months out of the year the rainfall is over 1000mm per month. Their primary crop productions are bananas and coffee; they use these as cash crops. Farmers still use substantial areas for millet farming. This is their main food crop. They have intensive home gardens to produce for the families' needs and the outlying farming is used mostly for cash crop production, where coffees and bananas are cultivated using intercropping methods.
Southern Africa
In Southern Africa, conglomerates of farming companies have primarily written-off the lands in Northwest Zambia. The land is mostly made of plateaus, with lower-lying lands in the Kabompo and Zambezi river valleys. Much of the area is in the Congo–Zambezi watershed, the area of land where all the sources of water run into the same river basin. The land has several types of soil. There are areas of yellow clay, with higher concentrations of sand, where higher rainfall causes soil erosion; this makes farming in those areas difficult. There are also areas with fertile red clay; these are rich soils good for cultivation of crops. Farmers in the area grow sorghum, millet, sweet potatoes, pumpkin, and maize. There is marginal maize production, used as a cash crop for the community. The area is usually covered by thick forestation that must be cleared before crops can be cultivated; they do not practice agroforestry. Traditionally, spouses must maintain their own gardens. The produce is shared communally, but the practice leads to intensive home gardens that are maintained by the women of the villages.
South Pacific
Highlander horticulture
The Enga of the Western Highlands Province in New Guinea receive most of their food from growing sweet potatoes Ipomoea batatas which they plant in mulch mounds at elevations up to 2,700 m or higher. The mounds that the Enga make to plant their crops of potatoes are formed from by piling large amounts of grass taken from fallow, or unplanted, plots then by covering the grass with dirt. The size of the mounds depends on elevation; the higher the elevation; the bigger the mounds will be. Mounds above 2,500 m in altitude can have a height of 0.85 m in height; while crops below 1,500 m are not mounded at all. The function of the mound is to protect the crops from the frequent frosts that occur at the high altitudes of the Enga. With sweet potatoes having a very long maturation period of nine months, the Enga also invest their time and space on the mounds with planting other crops that have much shorter maturation periods such as peas in case a heavy frost does claim the crop.
The planting of the mounds is done so that the plants which have a higher frost tolerance, such as the Irish potatoes, are planted casually throughout the mound and the low tolerant sweet potatoes are planted in the best position to avoid the frost. Peas, beans, and cabbage which are all highly tolerant to frost will be planted outside the circle of sweet potatoes and lower on the mound placing them closer to the cold temperatures of the ground. The Enga practice fallow rotation where a garden will in crop for about four years followed by about four years of fallow grassland to let the soil replenish.
Garden size for an average Enga garden is about 0.21 hectare or about 2,100 square meters and can contain a few hundred mounds. Another gardening strategy the Enga have implemented is the use of kin lands that are usually within one to two days walk from the farmers normal planting grounds. The uses of multiple gardens at differing elevation and the ability access clan lands in different areas for gardening have allowed the Enga to adapt to their environment and survive under harsh conditions.
Lowland Swidden cultivation
Swidden cultivation is an extensive agriculture practice that is also known as slash-and-burn agriculture. The process is extensive because it requires a vast amount of land divided into several plots with one plot planted for a period of years, while the other plots lay fallow for a number of years.
For the Bine-speaking peoples of the New Guinea lowlands swidden cultivation is a main practice for crop propagation. The main crop the Bine grow is the taro root, although they grow about 15 subsidiary crops including: sweet potato, banana, manioc, maize, yam, pawpaw, sugar cane, pineapple, and others. The swiddens which can be placed in either savannas or forests are created by cutting down all the vegetation in the area that the swidden will be. The farmers then pile all of the cut vegetation on the swidden plot and leave it to dry out through the dry season. Right before the wet season begins the piles are burned and the soil and ash are tilled together. The process of tilling the soil and ash mixes the carbon and nitrogen rich ash into the soil thereby fertilizing the soil for the coming crop. After the soil is tilled the crops are planted.
There are two planting years for a single swidden for the Bine farmers. In the first year the Bine plant primary taro root with a few subsidiary crops like bananas and sweet potatoes. In the second year taro root makes up about 50 percent of the swidden and the rest of the swidden is mixed with about 15 other plants. After the second year the Bine farmers move on to an adjacent swidden and allow the previous swidden to lay fallow or unplanted for a period of 5 to 10 years in order to repopulate the vegetation. The number of years that a swidden will lay fallow is determined by the plants demand for the nitrogen in the soil. Some plants will leach the soil of nitrogen in a few years and require four or five times that fallow; while other plants can be planted for many years and lay fallow only one or two times the planting period. Swidden cultivation requires a lot of land in order to feed only a few people, but the Bine, whose numbers are low, make good uses of their land through swidden farming.
Island horticulture
For most South Pacific Island cultures the main subsistence techniques are hunting and gathering. Fishing and the gathering of sago, banana, and other tropical foods are the norm with very little organized agriculture. The Tabalu of Kiriwina located in the Trobriand Islands practice a form of agriculture called Kaylu'ebila, a form of garden magic. The main crop for the Tabalu is the yam and there is a definite division of labour according to sex when it comes to gardening. Heavy work is done by the men and it includes clearing the vegetation, caring the yam supports, and planting the yam tubers in the ground. The women aid by weeding the gardens.
Gardening for the Tabalu is a very long and in depth magical process; with special magicians and magical ingredients which have been handed down from family member to family member over time. Garden fields which are called Kwabila are fenced in on all sides to keep out the swine that are breed by the Tabalu. Kwabila are then divided into many smaller plots called baleko, these are the individual gardens that the crops will be planted in.
South America
South America consists of modern-day Venezuela, Colombia, Ecuador, Suriname, Brazil, Peru, Bolivia, Paraguay, Uruguay, Argentina, and Chile. South America has historically been a land exploited not only for its natural resources, but also for its indigenous knowledge and labor force. The environmental diversification of South America has been at the foundation of its presence in the global economy as a resource for agriculture, forestry, fishing, hunting, livestock, mining and quarrying.
South America can be seen as cultural regions inhabited by marginal tribes, tropical forest tribes, and the circum-Caribbean tribes each with its own distinct way of agricultural cultivation. South America's geographic regions are inhabited by regional tribes which include; the Chocó in the Northern Columbian area, The Kayapo in the Eastern Para area, the Chono in the Southern Fuegian area, and the Quechuas in the Western Peruvian area, Each of these regions has adapted not only their own cultural identity and agricultural style,
The Eastern region of South America is known as the Para area in what is now Brazil, and has for millennia been home to the tropical forest Kayapớ tribe. The Kayapớ lived in sedentary villages and were proficient in pottery and loom-weaving, yet they did not domesticate animals or poses knowledge on metallurgy. These Tropical forest Tribes can be characterized through their farming, dugout canoes, woven baskets, loom weaving, and pole and thatch houses.
In the Para area the Kayapớ like most tribes of this region practiced intensive agriculture or clearing cultivation. Beginning their agricultural year with a low water season intensifies fishing. The low water season is then followed by the high water season or harvest season. It is during this harvest season that the Kayapớ are able to exercise their leisure time before the cycle ends with (low water levels) and a return to intense fishing. Each changing season commences ceremonies for the Kayapớ that are directly tied to agriculture, hunting, or fishing. Unlike the Chocó the Kayapớ used an agricultural method known as the slash-burn method (shifting agriculture). The Kayapớ cut the forest in April to September (dry season) and time their burns just before the raining season. The Kayapớ used circular plots for agricultural cultivation consisting of five rings (characterized as cultivation zones). The first circle or the inner circle was used for taro and sweet potatoes that thrive in the hotter soil found in the center of the plot. The second circle cultivated maize, manioc, and rice an area that needed various ash enrichment treatments and would experience short term growth. The third zone was an area of rich soil and best served mixed crops including the banana, urucu, papaya. The fourth zones consists of shade loving plants and are for a medical purpose, yet evidence of beans and yams have been also found here. The fifth zone or the outer ring was left as a protective zone that included trap animals protective insects and birds. This form of agriculture requires not only intense physical labor but also requires a knowledge of not only the land, but various types of ground cover, shades and temperatures of local soils, as well as cloud formations to time careful burning. When the Kayapớ manage their agricultural plots they must work with a variety of interacting factors including the background soil fertility, the heterogeneous quality of ash and its distribution, crop nutrient requirements, cropping cycles, management requirements, and pest and disease control clearly illustrating the common misconception that this form of agriculture is primitive and ineffiecnet.
It has often been thought that the slash-burn plots are abandoned after one or two years because of un productive soil, but this is a common misconception. The Kayapớ revisit abandoned fields because plants can offer direct and indirect benefits. One direct benefit would be the ability to eat that which has been produced and an indirect benefit would be that open fields allow attract game for hunting and can produce long after they have been tended.
The Southern region of South America is known as the Fuegian area and is occupied by the Chono, Alacaluf and Yahgan. These Marginal tribes differed greatly from the other regions in that they were expert in making bark or plank canoes and domesticating dogs, hunting, fishing and gathering. Nomads with simple socio-religious patterns, yet completely lacked the technology of pottery, loom-weaving, metallurgy and even agriculture. Since there was a lack of agriculture in this region the Chono ate native berries, roots, fruits, and wild celery. The source of nutrition for the Chono, Yahgans, and Alcaaluf thus mainly consisted of sea food such as; whales, seal porpoises, guanacos, and otters. The Southern tribes of South America distinguished themselves not in having a rare form of agriculture like the North's slush-mush method or the having intensive agriculture like the East's sophisticated slash-burn method, yet were able to distinguished themselves in being absent of this trait.
North America
Farming methods developed by Native Americans include terracing, irrigation, mound building, crop rotation and fertilization. They also used extensive companion planting (see the Three Sisters).
Terracing is an effective technique in a steep-sloped, semi-arid climate. The Indigenous farmers stair-stepped the hills so that soil erosion was minimal and land surface was better suited for farming. In the Southwest, including parts of New Mexico, Arizona, and parts of Northern Mexico, terracing was extensive. Terraces were constructed by placing rock dams to redirect runoff water to canals that evenly dispersed rain water. The terraced field transformed the terrain into land suitable for farming maize. There is evidence that terracing has been used in the Southwest for about 2,500 years. The Anasazi people from this region built reservoirs and directed rain water through ditches to water the crops in the terraces. The natives grew corn, squash, and beans, along with other crops in the terraced fields.
Corn, squash, and beans were staple crops for Native Americans and were grown throughout much of the North American continent. This trio is known as the Three sisters. Ancient folklore belief says that the Three Sisters represented three goddesses. Each sister protected the other two, and therefore the Three Sisters should never be separated and instead be planted, cooked, and consumed together. In reality, This Triad was an example of symbiotic planting. The corn stalks functioned as a support for the beans. The beans fixed nitrogen into a usable form for the corn and squash, and the broad squash leaves provided shade for the soil, which aided in preventing evaporation and controlling weeds. As the success of the Three Sisters spread, many cultures turned away from hunting and gathering and relied much more on farming. Geographically native cultures in the Woodlands, Prairie, Plains, Great Basin, and Plateau regions of North America all utilized the Three Sisters to some extent. Where they were not grown, the locals traded for them. Nomadic tribes, such as the Dakotas, would trade meat for these vegetable staples. The Three Sisters were usually eaten in unison, as they provide fairly balanced nutrition when consumed together(for example, beans and corn together provide a complete set of the essential amino acids).
Native American farmers also employed irrigation. This technique was utilized throughout much of the Southwest and is useful where water is scarce. Irrigation was and is still used today throughout much of the world. Native Americans controlled the amount of water that reached their fields by building long irrigation canals to redirect water from a source to water their crops. The Hohokam people constructed about 600 miles of irrigation canals from AD 50 to 1450 near Phoenix, Arizona. Part of the canal is used by the City of Phoenix today. The Olmecs of Mesoamerica built canals over 4,000 years ago.
Chinampas, artificial islands constructed in swamplands and lakes, were invented by Mayan farmers and the technique became used extensively throughout Mesoamerica and was later used by the Aztecs as part of the land reclamation process of the city of Tenochtitlan. This technique increases arable land, and provided additional farming plots the population of Mesoamerica grew.
In the Northeast Woodlands and the Great Lakes region, an advanced society known as the Moundbuilders emerged. This society lived in the flood plains of the Mississippi river basin. This culture farmed mainly maize. They had little need for foraging and grew to an advanced civilization due to food surplus. They were the largest civilization north of the Rio Grande.
Native Americans also developed storage systems such as storage containers which allowed them to store seeds to plant during the next planting season. They also stored food in dug-out pits or holes in hillsides. Native Americans developed corn cribs. These were storage bins that were elevated off the ground. This technique prevented moisture and animal intrusion.
Selective crop breeding was also employed. Corn is a domestic plant and cannot grow on its own. The first corn grown by Native Americans had small ears, and only produced a few kernels per ear. By 2,000 years ago, single stalks with large ears were being produced. Native Americans created over 700 varieties of corn by 1500 AD.
"Hands-off" hunter and gatherer misconception
Anderson breaks down common misconceptions surrounding the way in which Native Americans lived among nature and shows that their true impact was often one that attracted and invited others to the land they inhabited. Native Americans, regardless of location, were privy to wildlife management such as, "coppicing, pruning, harrowing, sowing, weeding, burning, digging, thinning, and selective harvesting." These practices were acted out in calculated intervals according to season and by witnessing the signals conveyed by nature to do so and, "on the whole, allowed for sustainable harvest of plants over centuries." Anderson's writing focuses specifically on Natives established in California and emphasizes these techniques applied by them to be essential in maintaining, and even creating, the rich Californian landscapes settlers later happened upon, such as: the coastal prairies, valley grasslands, and oak savannas.
Overall, the Native's practices were nature conserving and sustaining due to the specificity of their environmental knowledge which was learned through more than twelve-thousand years of trial and error. Anderson also establishes that the Native's practices were most likely not always beneficial or environmentally friendly. Though he asserts that there is no real evidence of their destruction available, it is possible that the Indigenous peoples in California may have been responsible for the extinction of early regional species.
In discussing California Native's "tempered" land tenure practices, Anderson deconstructs the idea of Native Americans as hunter-gatherers whose sparse population and nomadic ways left little to no mark on the land they traveled. When European and Asian farmers, ranchers, and entrepreneurs established themselves in that same land, the concept of their surroundings that they perpetuated was one of an uninhabited wilderness, untouched by man. Though as Anderson points out, this was never the case. This "wilderness" was carefully tended to by the Natives for hundreds of years who had both negative and positive impacts on its conservation. It is theorized by Anderson that without the Natives' calculated intervention, real wilderness in the form of thickets, dense understory, and wildfire would have deemed the land uninhabitable. In establishing this belief, the foreign settlers not only erased the Native's long history of masterful cultivation from the land, but they pushed their people out of the land as well, establishing a new, Eurocentric construction of the American's continent. This constructed history has not only diminished Native ancestry and culture, but also the land, which is no longer maintained or protected by strategic resource management. The European way of mitigating resource and land depletion follows a "hands off" model, which comes from the misconceptions of "leave no trace" mentioned previously.
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22661773 | https://en.wikipedia.org/wiki/Animal-free%20agriculture | Animal-free agriculture | Animal-free agriculture, also known as plant agriculture, plant-based agriculture, veganic agriculture, stockfree farming, plant farming or veganic farming, consists of farming methods that do not use animals or animal products.
Animal-free growers do not keep domesticated animals and do not use animal products such as farmed animal manures or animal parts (bone meal, blood meal, fish meal) to fertilize their crops. Emphasis is placed on using green manures and plant-based compost instead.
Methods
Animal-free farming may use organic or non-organic farming methods. However, most detailed discussions of animal-free agriculture currently focus on animal-free organic variants. In the European Union, farmers have a financial incentive to use manure instead of animal-free fertilisers, since manure is subsidised. However, organic manure is not subsidised.
Industrial agriculture with synthetic fertilizers is animal-free. In the United States, few industrial farms use manure. Of all U.S. cropland, only 5% was manured in 2006.
Vegan organic farming
Vegan organic farming methods do not use animal products or by-products, such as blood meal, fish products, bone meal, feces, or other animal-origin matter because the production of these materials is viewed as either harming animals directly, or as associated with the exploitation and consequent suffering of animals. Some of these materials are by-products of animal husbandry, created during the process of cultivating animals for the production of meat, milk, skins, furs, entertainment, labor, or companionship. The sale of such by-products decreases expenses and increases profit for those engaged in animal husbandry and therefore helps support the animal husbandry industry, an outcome vegans find unacceptable.
Vegan organic growers maintain soil fertility using green manures, cover crops, green wastes, composted vegetable matter, and minerals. Some vegan gardeners may supplement this with human urine from vegans (which provides nitrogen) and 'humanure' from vegans, produced from compost toilets.
Veganic organic farmers take measures such as refraining from making large disturbances in the soil of the land and cultivating a variety of plants in the ground. This form of farming "encompasses a respect for the animals, the environment, and human health." Some of the plant-based techniques used in veganic agriculture include mulch, compost, chipped branched wood, crop rotation and others.
Farms certified as biocyclic vegan use preventative methods to manage insects. If these fail, however, the label allows them to use insecticides such as Bacillus thuringiensis, which starves larvae to death.
Vegan organic farming is much less common than organic farming. In 2019, there were 63 self-declared vegan organic farms in the United States, and 16,585 certified organic farms.
Timeline
2006
The World Conservation Union's Red List of Threatened Species reports that most of the world's threatened species are experiencing habitat loss as a result of livestock production conducted through animal agriculture.
Center for Science in the Public Interest releases Six Arguments for a Greener Diet that found that a plant-rich diet "leads to much less food poisoning, water pollution, air pollution, global warming."
2016
Research published in the journal Nature Communications finds that vegan diets have the best land use and are the only way to feed the global population by 2050.
The World Resources Institute published the report: Shifting Diets for a Sustainable Food Future which showed that if people who consume large amounts of meat and dairy changed to diets with more plant-based meals could reduce agriculture's pressure on the environment.
2017
University of Edinburgh researchers find that animal farming is the leading cause of food waste as it is responsible for the most losses of all harvested crops on Earth (40%) due to secondary consumption.
Forbes magazine publishes a compilation of recent vegan and plant-based business successes noting that vegan living is becoming more a norm because of its positive impact on sustainability.
2018
Research published in the Proceedings of the National Academy of Sciences find that a vegan shift would increase the US food supply by a third, eliminating all of the losses due to food waste and feeding all Americans as well as roughly 390,000,000 more.
A Harvard study found that shifting all beef production in the U.S. to pastured, grass-fed systems would require 30% more cattle, increase beef's methane emissions by 43%, and would require much more land than is available.
2019
A report from the Humane Party determines that vegan-organic agriculture can be 4,198% more productive than animal-based agriculture in the amount of food produced per acre.
Veganic farmer Will Bonsall told The Guardian that most vegetables are "very un-vegan" due to being grown using inputs of animal-based products.
Advantages
Livestock in the United States produce 230,000 pounds of manure per second, and nitrogen from these wastes is converted into ammonia and nitrates which leach into ground and surface water causing contamination of wells, rivers and streams. Mature compost of plant-based origins, used in animal-free agriculture, can reduce leaching of nitrate which leads to an improvement in groundwater quality and counteracts the eutrophication of surface waters.
Animal free agriculture has the potential to prevent illnesses like influenza from spreading. Experts agree that most strains of the influenza virus that infect human beings came from contact with other animals. Farm animals on factory farms may be genetically similar therefore making them more susceptible to specific parasites. Infection among animals is more easily spread because of their close proximity to one another. Animal-free agriculture does not contribute to the spread of influenza through animals.
Current use
Vegan France Interpro in collaboration with the Biocyclic Vegan Network created an interactive map that lists all-vegan organic projects across Europe. This list primarily includes agricultural operations but also trading and processing companies, online shops, network organizations as well as certification bodies that certify farms according to the Biocyclic Vegan Standard.
There is a similar map in North America that conducts the same concept and locates vegan farms around North America.
The Biocyclic Vegan Standard is an IFOAM-accredited organic standard for vegan organic farms. It is awarded by BNS Biocyclic Network Services Ltd (a Cypriot company), and has accredited 19 farms in Europe . The German Environment Agency awarded the German biocyclic vegan association some 60,000 euros for the promotion of the biocyclic vegan standard from 2021 to 2022.
, 18 farms in the United Kingdom and Ireland are certified vegan organic by the Stockfree Organic label. Farms wanting to obtain the label are certified by the Soil Association, and the label's requirements are determined by the Vegan Organic Network.
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20082214 | https://en.wikipedia.org/wiki/Obsessive%E2%80%93compulsive%20disorder | Obsessive–compulsive disorder | Obsessive–compulsive disorder (OCD) is a mental and behavioral disorder in which an individual has intrusive thoughts (an obsession) and feels the need to perform certain routines (compulsions) repeatedly to relieve the distress caused by the obsession, to the extent where it impairs general function.
Obsessions are persistent unwanted thoughts, mental images or urges that generate feelings of anxiety, disgust or discomfort. Some common obsessions include fear of contamination, obsession with symmetry, the fear of acting blasphemously, the sufferer's sexual orientation and the fear of possibly harming others or themselves. Compulsions are repeated actions or routines that occur in response to obsessions to achieve a relief from anxiety. Common compulsions include excessive hand washing, cleaning, counting, ordering, repeating, avoiding triggers, hoarding, neutralizing, seeking assurance, praying and checking things. People with OCD may only perform mental compulsions such as needing to know or remember things. While this is sometimes referred to as primarily obsessional obsessive–compulsive disorder (Pure O), it is also considered a misnomer due to associated mental compulsions and reassurance seeking behaviors that are consistent with OCD.
Compulsions occur often and typically take up at least one hour per day, impairing one's quality of life. Compulsions cause relief in the moment, but cause obsessions to grow over time due to the repeated reward-seeking behavior of completing the ritual for relief. Many adults with OCD are aware that their compulsions do not make sense, but they still perform them to relieve the distress caused by obsessions. For this reason, thoughts and behaviors in OCD are usually considered egodystonic (inconsistent with one's ideal self-image). In contrast, thoughts and behaviors in obsessive–compulsive personality disorder (OCPD) are usually considered egosyntonic (consistent with one's ideal self-image), helping differentiate between OCPD and OCD.
Although the exact cause of OCD is unknown, several regions of the brain have been implicated in its neuroanatomical model including the anterior cingulate cortex, orbitofrontal cortex, amygdala and BNST. The presence of a genetic component is evidenced by the increased likelihood for both identical twins to be affected than both fraternal twins. Risk factors include a history of child abuse or other stress-inducing events such as during the postpartum period or after streptococcal infections. Diagnosis is based on clinical presentation and requires ruling out other drug-related or medical causes; rating scales such as the Yale–Brown Obsessive–Compulsive Scale (Y-BOCS) assess severity. Other disorders with similar symptoms include generalized anxiety disorder, major depressive disorder, eating disorders, tic disorders, body-focused repetitive behavior and obsessive–compulsive personality disorder. Personality disorders are a common comorbidity, with schizotypal and OCPD having poor treatment response. The condition is also associated with a general increase in suicidality. The phrase obsessive–compulsive is sometimes used in an informal manner unrelated to OCD to describe someone as excessively meticulous, perfectionistic, absorbed or otherwise fixated. However, the actual disorder can vary in presentation and individuals with OCD may not be concerned with cleanliness or symmetry.
OCD is chronic and long-lasting with periods of severe symptoms followed by periods of improvement. Treatment can improve ability to function and quality of life, and is usually reflected by improved Y-BOCS scores. Treatment for OCD may involve psychotherapy, pharmacotherapy such as antidepressants or surgical procedures such as deep brain stimulation or, in extreme cases, psychosurgery. Psychotherapies derived from cognitive behavioral therapy (CBT) models, such as exposure and response prevention, acceptance and commitment therapy, and inference based-therapy, are more effective than non-CBT interventions. Selective serotonin reuptake inhibitors (SSRIs) are more effective when used in excess of the recommended depression dosage; however, higher doses can increase side effect intensity. Commonly used SSRIs include sertraline, fluoxetine, fluvoxamine, paroxetine, citalopram and escitalopram. Some patients fail to improve after taking the maximum tolerated dose of multiple SSRIs for at least two months; these cases qualify as treatment-resistant and can require second-line treatment such as clomipramine or atypical antipsychotic augmentation. While SSRIs continue to be first-line, recent data for treatment-resistant OCD supports adjunctive use of neuroleptic medications, deep brain stimulation and neurosurgical ablation. There is growing evidence to support the use of deep brain stimulation and repetitive transcranial magnetic stimulation for treatment-resistant OCD.
Obsessive–compulsive disorder affects about 2.3% of people at some point in their lives, while rates during any given year are about 1.2%. More than three million Americans suffer from OCD. According to Mercy, approximately 1 in 40 U.S. adults and 1 in 100 U.S. children have OCD. Although possible at times with triggers such as pregnancy, onset rarely occurs after age 35 and about 50% of patients experience detrimental effects to daily life before age 20. While OCD occurs worldwide, a recent meta-analysis showed that women are 1.6 times more likely to experience OCD. Based on data from 34 studies, the worldwide prevalence rate is 1.5% in women and 1% in men.
Signs and symptoms
OCD can present with a wide variety of symptoms. Certain groups of symptoms usually occur together as dimensions or clusters, which may reflect an underlying process. The standard assessment tool for OCD, the Yale–Brown Obsessive Compulsive Scale (Y-BOCS), has 13 predefined categories of symptoms. These symptoms fit into three to five groupings. A meta-analytic review of symptom structures found a four-factor grouping structure to be most reliable: symmetry factor, forbidden thoughts factor, cleaning factor and hoarding factor. The symmetry factor correlates highly with obsessions related to ordering, counting and symmetry, as well as repeating compulsions. The forbidden thoughts factor correlates highly with intrusive thoughts of a violent, religious or sexual nature. The cleaning factor correlates highly with obsessions about contamination and compulsions related to cleaning. The hoarding factor only involves hoarding-related obsessions and compulsions, and was identified as being distinct from other symptom groupings.
When looking into the onset of OCD, one study suggests that there are differences in the age of onset between males and females, with the average age of onset of OCD being 9.6 for male children and 11.0 for female children. Children with OCD often have other mental disorders, such as ADHD, depression, anxiety and disruptive behavior disorder. Continually, children are more likely to struggle in school and experience difficulties in social situations (Lack 2012). When looking at both adults and children a study found the average ages of onset to be 21 and 24 for males and females respectively. While some studies have shown that OCD with earlier onset is associated with greater severity, other studies have not been able to validate this finding. Looking at women specifically, a different study suggested that 62% of participants found that their symptoms worsened at a premenstrual age. Across the board, all demographics and studies showed a mean age of onset of less than 25.
Some OCD subtypes have been associated with improvement in performance on certain tasks, such as pattern recognition (washing subtype) and spatial working memory (obsessive thought subtype). Subgroups have also been distinguished by neuroimaging findings and treatment response, though neuroimaging studies have not been comprehensive enough to draw conclusions. Subtype-dependent treatment response has been studied and the hoarding subtype has consistently been least responsive to treatment.
While OCD is considered a homogeneous disorder from a neuropsychological perspective, many of the symptoms may be the result of comorbid disorders. For example, adults with OCD have exhibited more symptoms of attention deficit hyperactivity disorder (ADHD) and autism spectrum disorder (ASD) than adults without OCD.
In regards to the cause of onset, researchers asked participants in one study what they felt was responsible for triggering the initial onset of their illness. 29% of patients answered that there was an environmental factor in their life that did so. Specifically, the majority of participants who answered with that noted their environmental factor to be related to an increased responsibility.
Obsessions
Obsessions are stress-inducing thoughts that recur and persist, despite efforts to ignore or confront them. People with OCD frequently perform tasks, or compulsions, to seek relief from obsession-related anxiety. Within and among individuals, initial obsessions vary in clarity and vividness. A relatively vague obsession could involve a general sense of disarray or tension, accompanied by a belief that life cannot proceed as normal while the imbalance remains. A more intense obsession could be a preoccupation with the thought or image of a close family member or friend dying, or intrusive thoughts related to relationship rightness. Other obsessions concern the possibility that someone or something other than oneself—such as God, the devil or disease—will harm either the patient or the people or things the patient cares about. Others with OCD may experience the sensation of invisible protrusions emanating from their bodies or feel that inanimate objects are ensouled. Another common obsession is scrupulosity, the pathological guilt/anxiety about moral or religious issues. In scrupulosity, a person's obsessions focus on moral or religious fears, such as the fear of being an evil person or the fear of divine retribution for sin. Mysophobia, a pathological fear of contamination and germs, is another common obsession theme.
Some people with OCD experience sexual obsessions that may involve intrusive thoughts or images of "kissing, touching, fondling, oral sex, anal sex, intercourse, incest and rape" with "strangers, acquaintances, parents, children, family members, friends, coworkers, animals and religious figures" and can include heterosexual or homosexual contact with people of any age. Similar to other intrusive thoughts or images, some disquieting sexual thoughts are normal at times, but people with OCD may attach extraordinary significance to such thoughts. For example, obsessive fears about sexual orientation can appear to the affected individual, and even to those around them, as a crisis of sexual identity. Furthermore, the doubt that accompanies OCD leads to uncertainty regarding whether one might act on the troubling thoughts, resulting in self-criticism or self-loathing.
Most people with OCD understand that their thoughts do not correspond with reality; however, they feel that they must act as though these ideas are correct or realistic. For example, someone who engages in compulsive hoarding might be inclined to treat inorganic matter as if it had the sentience or rights of living organisms, despite accepting that such behavior is irrational on an intellectual level. There is debate as to whether hoarding should be considered an independent syndrome from OCD.
Compulsions
Some people with OCD perform compulsive rituals because they inexplicably feel that they must do so, while others act compulsively to mitigate the anxiety that stems from obsessive thoughts. The affected individual might feel that these actions will either prevent a dreaded event from occurring or push the event from their thoughts. In any case, their reasoning is so idiosyncratic or distorted that it results in significant distress, either personally or for those around the affected individual. Excessive skin picking, hair pulling, nail biting and other body-focused repetitive behavior disorders are all on the obsessive–compulsive spectrum. Some individuals with OCD are aware that their behaviors are not rational, but they feel compelled to follow through with them to fend off feelings of panic or dread. Furthermore, compulsions often stem from memory distrust, a symptom of OCD characterized by insecurity in one's skills in perception, attention and memory, even in cases where there is no clear evidence of a deficit.
Common compulsions may include hand washing, cleaning, checking things (such as locks on doors), repeating actions (such as repeatedly turning on and off switches), ordering items in a certain way and requesting reassurance. Although some individuals perform actions repeatedly, they do not necessarily perform these actions compulsively; for example, morning or nighttime routines and religious practices are not usually compulsions. Whether behaviors qualify as compulsions or mere habit depends on the context in which they are performed. For instance, arranging and ordering books for eight hours a day would be expected of someone who works in a library, but this routine would seem abnormal in other situations. In other words, habits tend to bring efficiency to one's life, while compulsions tend to disrupt it. Furthermore, compulsions are different from tics (such as touching, tapping, rubbing or blinking) and stereotyped movements (such as head banging, body rocking or self-biting), which are usually not as complex and not precipitated by obsessions. It can sometimes be difficult to tell the difference between compulsions and complex tics, and about 10–40% of people with OCD also have a lifetime tic disorder.
People with OCD rely on compulsions as an escape from their obsessive thoughts; however, they are aware that relief is only temporary and that intrusive thoughts will return. Some affected individuals use compulsions to avoid situations that may trigger obsessions. Compulsions may be actions directly related to the obsession, such as someone obsessed with contamination compulsively washing their hands, but they can be unrelated as well. In addition to experiencing the anxiety and fear that typically accompanies OCD, affected individuals may spend hours performing compulsions every day. In such situations, it can become difficult for the person to fulfill their work, familial or social roles. These behaviors can also cause adverse physical symptoms; for example, people who obsessively wash their hands with antibacterial soap and hot water can make their skin red and raw with dermatitis.
Individuals with OCD often use rationalizations to explain their behavior; however, these rationalizations do not apply to the behavioral pattern, but to each individual occurrence. For example, someone compulsively checking the front door may argue that the time and stress associated with one check is less than the time and stress associated with being robbed, and checking is consequently the better option. This reasoning often occurs in a cyclical manner and can continue for as long as the affected person needs it to in order to feel safe.
In cognitive behavioral therapy (CBT), OCD patients are asked to overcome intrusive thoughts by not indulging in any compulsions. They are taught that rituals keep OCD strong, while not performing them causes OCD to become weaker. This position is supported by the pattern of memory distrust; the more often compulsions are repeated, the more weakened memory trust becomes and this cycle continues as memory distrust increases compulsion frequency. For body-focused repetitive behaviors (BFRB) such as trichotillomania (hair pulling), skin picking and onychophagia (nail biting), behavioral interventions such as habit reversal training and decoupling are recommended for the treatment of compulsive behaviors.
OCD sometimes manifests without overt compulsions, which may be termed "primarily obsessional OCD." OCD without overt compulsions could, by one estimate, characterize as many as 50–60% of OCD cases.
Insight and overvalued ideation
The Diagnostic and Statistical Manual of Mental Disorders (DSM-5), identifies a continuum for the level of insight in OCD, ranging from good insight (the least severe) to no insight (the most severe). Good or fair insight is characterized by the acknowledgment that obsessive–compulsive beliefs are not or may not be true, while poor insight, in the middle of the continuum, is characterized by the belief that obsessive–compulsive beliefs are probably true. The absence of insight altogether, in which the individual is completely convinced that their beliefs are true, is also identified as a delusional thought pattern and occurs in about 4% of people with OCD. When cases of OCD with no insight become severe, affected individuals have an unshakable belief in the reality of their delusions, which can make their cases difficult to differentiate from psychotic disorders.
Some people with OCD exhibit what is known as overvalued ideas, ideas that are abnormal compared to affected individuals' respective cultures, and more treatment-resistant than most negative thoughts and obsessions. After some discussion, it is possible to convince the individual that their fears are unfounded. It may be more difficult to practice exposure and response prevention therapy (ERP) on such people, as they may be unwilling to cooperate, at least initially. Similar to how insight is identified on a continuum, obsessive-compulsive beliefs are characterized on a spectrum, ranging from obsessive doubt to delusional conviction. In the United States, overvalued ideation (OVI) is considered most akin to poor insight—especially when considering belief strength as one of an idea's key identifiers. Furthermore, severe and frequent overvalued ideas are considered similar to idealized values, which are so rigidly held by, and so important to affected individuals, that they end up becoming a defining identity. In adolescent OCD patients, OVI is considered a severe symptom.
Historically, OVI has been thought to be linked to poorer treatment outcome in patients with OCD, but it is currently considered a poor indicator of prognosis. The Overvalued Ideas Scale (OVIS) has been developed as a reliable quantitative method of measuring levels of OVI in patients with OCD and research has suggested that overvalued ideas are more stable for those with more extreme OVIS scores.
Cognitive performance
Though OCD was once believed to be associated with above-average intelligence, this does not appear to necessarily be the case. A 2013 review reported that people with OCD may sometimes have mild but wide-ranging cognitive deficits, most significantly those affecting spatial memory and to a lesser extent with verbal memory, fluency, executive function and processing speed, while auditory attention was not significantly affected. People with OCD show impairment in formulating an organizational strategy for coding information, set-shifting, and motor and cognitive inhibition.
Specific subtypes of symptom dimensions in OCD have been associated with specific cognitive deficits. For example, the results of one meta-analysis comparing washing and checking symptoms reported that washers outperformed checkers on eight out of ten cognitive tests. The symptom dimension of contamination and cleaning may be associated with higher scores on tests of inhibition and verbal memory.
Video game addiction
Pediatric OCD
Approximately 1–2% of children are affected by OCD. There is a lot of similarity between the clinical presentation of OCD in children and adults and it is considered a highly familial disorder, with a phenotypic heritability of around 50%. Obsessive–compulsive disorder symptoms tend to develop more frequently in children 10–14 years of age, with males displaying symptoms at an earlier age, and at a more severe level than females. In children, symptoms can be grouped into at least four types, including sporadic and tic-related OCD.
The Children's Yale–Brown Obsessive–Compulsive Scale (CY-BOCS) is the gold standard measure for assessment of pediatric OCD. It follows the Y-BOCS format, but with a Symptom Checklist that is adapted for developmental appropriateness. Insight, avoidance, indecisiveness, responsibility, pervasive slowness and doubting are not included in a rating of overall severity. The CY-BOCS has demonstrated good convergent validity with clinician-rated OCD severity and good to fair discriminant validity from measures of closely related anxiety, depression and tic severity. The CY-BOCS Total Severity score is an important monitoring tool as it is responsive to pharmacotherapy and psychotherapy. Positive treatment response is characterized by 25% reduction in CY-BOCS total score and diagnostic remission is associated with a 45%-50% reduction in Total Severity score (or a score <15).
CBT is the first line treatment for mild to moderate cases of OCD in children, while medication plus CBT is recommended for moderate to severe cases. Serotonin reuptake inhibitors (SRIs) are first-line medications for OCD in children with established AACAP guidelines for dosing.
Associated conditions
People with OCD may be diagnosed with other conditions as well, such as obsessive–compulsive personality disorder, major depressive disorder, bipolar disorder, generalized anxiety disorder, anorexia nervosa, social anxiety disorder, bulimia nervosa, Tourette syndrome, transformation obsession, ASD, ADHD, dermatillomania, body dysmorphic disorder and trichotillomania. More than 50% of people with OCD experience suicidal tendencies and 15% have attempted suicide. Depression, anxiety and prior suicide attempts increase the risk of future suicide attempts.
It has been found that between 18 and 34% of females currently experiencing OCD scored positively on an inventory measuring disordered eating. Another study found that 7% are likely to have an eating disorder, while another found that fewer than 5% of males have OCD and an eating disorder.
Individuals with OCD have also been found to be affected by delayed sleep phase disorder at a substantially higher rate than the general public. Moreover, severe OCD symptoms are consistently associated with greater sleep disturbance. Reduced total sleep time and sleep efficiency have been observed in people with OCD, with delayed sleep onset and offset.
Some research has demonstrated a link between drug addiction and OCD. For example, there is a higher risk of drug addiction among those with any anxiety disorder, likely as a way of coping with the heightened levels of anxiety. However, drug addiction among people with OCD may be a compulsive behavior. Depression is also extremely prevalent among people with OCD. One explanation for the high depression rate among OCD populations was posited by Mineka, Watson and Clark (1998), who explained that people with OCD, or any other anxiety disorder, may feel "out of control".
Someone exhibiting OCD signs does not necessarily have OCD. Behaviors that present as obsessive–compulsive can also be found in a number of other conditions, including obsessive–compulsive personality disorder (OCPD), autism spectrum disorder (ASD) or disorders in which perseveration is a possible feature (ADHD, PTSD, bodily disorders or stereotyped behaviors). Some cases of OCD present symptoms typically associated with Tourette syndrome, such as compulsions that may appear to resemble motor tics; this has been termed tic-related OCD or Tourettic OCD.
OCD frequently occurs comorbidly with both bipolar disorder and major depressive disorder. Between 60 and 80% of those with OCD experience a major depressive episode in their lifetime. Comorbidity rates have been reported at between 19 and 90%, as a result of methodological differences. Between 9–35% of those with bipolar disorder also have OCD, compared to 1–2% in the general population. About 50% of those with OCD experience cyclothymic traits or hypomanic episodes. OCD is also associated with anxiety disorders. Lifetime comorbidity for OCD has been reported at 22% for specific phobia, 18% for social anxiety disorder, 12% for panic disorder and 30% for generalized anxiety disorder. The comorbidity rate for OCD and ADHD has been reported to be as high as 51%.
Causes
The cause of OCD is unknown. Both environmental and genetic factors are believed to play a role. Risk factors include a history of adverse childhood experiences or other stress-inducing events.
Drug-induced OCD
Some medications, toxin exposures and drugs, such as methamphetamine or cocaine, can induce obsessive–compulsive symptoms in people without a history of OCD. Atypical antipsychotics such as olanzapine and clozapine can induce OCD in some people, particularly individuals with schizophrenia.
The diagnostic criteria include:
General OCD symptoms (obsessions, compulsions, skin picking, hair pulling, etc.) that developed soon after exposure to the substance or medication which can produce such symptoms.
The onset of symptoms cannot be explained by an obsessive–compulsive and related disorder that is not substance/medication-induced and should last for a substantial period of time (about 1 month)
This disturbance does not only occur during delirium.
Clinically induces distress or impairment in social, occupational or other important areas of functioning.
Genetics
There appear to be some genetic components of OCD causation, with identical twins more often affected than fraternal twins. Furthermore, individuals with OCD are more likely to have first-degree family members exhibiting the same disorders than matched controls. In cases in which OCD develops during childhood, there is a much stronger familial link in the disorder than with cases in which OCD develops later in adulthood. In general, genetic factors account for 45–65% of the variability in OCD symptoms in children diagnosed with the disorder. A 2007 study found evidence supporting the possibility of a heritable risk for OCD. OCD is believed to be a heterogeneous disorder.
Research has found there to be a genetic correlation between anorexia nervosa and OCD, suggesting a strong etiology. First and second hand relatives of probands with OCD have a greater risk of developing anorexia nervosa as genetic relatedness increases.
A mutation has been found in the human serotonin transporter gene hSERT in unrelated families with OCD.
A systematic review found that while neither allele was associated with OCD overall, in Caucasians, the L allele was associated with OCD. Another meta-analysis observed an increased risk in those with the homozygous S allele, but found the LS genotype to be inversely associated with OCD.
A genome-wide association study found OCD to be linked with single-nucleotide polymorphisms (SNPs) near BTBD3 and two SNPs in DLGAP1 in a trio-based analysis, but no SNP reached significance when analyzed with case-control data.
One meta-analysis found a small but significant association between a polymorphism in SLC1A1 and OCD.
The relationship between OCD and Catechol-O-methyltransferase (COMT) has been inconsistent, with one meta-analysis reporting a significant association, albeit only in men, and another meta analysis reporting no association.
It has been postulated by evolutionary psychologists that moderate versions of compulsive behavior may have had evolutionary advantages. Examples would be moderate constant checking of hygiene, the hearth or the environment for enemies. Similarly, hoarding may have had evolutionary advantages. In this view, OCD may be the extreme statistical tail of such behaviors, possibly the result of a high number of predisposing genes.
Brain structure and functioning
Imaging studies have shown differences in the frontal cortex and subcortical structures of the brain in patients with OCD. There appears to be a connection between the OCD symptoms and abnormalities in certain areas of the brain, but such a connection is not clear. Some people with OCD have areas of unusually high activity in their brain or low levels of the chemical serotonin, which is a neurotransmitter that some nerve cells use to communicate with each other, and is thought to be involved in regulating many functions, influencing emotions, mood, memory and sleep.
Autoimmune
A controversial hypothesis is that some cases of rapid onset of OCD in children and adolescents may be caused by a syndrome connected to Group A streptococcal infections (GABHS), known as pediatric autoimmune neuropsychiatric disorders associated with streptococcal infections (PANDAS). OCD and tic disorders are hypothesized to arise in a subset of children as a result of a post-streptococcal autoimmune process. The PANDAS hypothesis is unconfirmed and unsupported by data and two new categories have been proposed: PANS (pediatric acute-onset neuropsychiatric syndrome) and CANS (childhood acute neuropsychiatric syndrome). The CANS and PANS hypotheses include different possible mechanisms underlying acute-onset neuropsychiatric conditions, but do not exclude GABHS infections as a cause in a subset of individuals. PANDAS, PANS and CANS are the focus of clinical and laboratory research, but remain unproven. Whether PANDAS is a distinct entity differing from other cases of tic disorders or OCD is debated.
A review of studies examining anti-basal ganglia antibodies in OCD found an increased risk of having anti-basal ganglia antibodies in those with OCD versus the general population.
Environment
OCD may be more common in people who have been bullied, abused or neglected, and it sometimes starts after a significant life event, such as childbirth or bereavement. It has been reported in some studies that there is a connection between childhood trauma and obsessive-compulsive symptoms. More research is needed to understand this relationship better.
Mechanisms
Neuroimaging
Functional neuroimaging during symptom provocation has observed abnormal activity in the orbitofrontal cortex (OFC), left dorsolateral prefrontal cortex (dlPFC), right premotor cortex, left superior temporal gyrus, globus pallidus externus, hippocampus and right uncus. Weaker foci of abnormal activity were found in the left caudate, posterior cingulate cortex and superior parietal lobule. However, an older meta-analysis of functional neuroimaging in OCD reported that the only consistent functional neuroimaging finding was increased activity in the orbital gyrus and head of the caudate nucleus, while anterior cingulate cortex (ACC) activation abnormalities were too inconsistent.
A meta-analysis comparing affective and nonaffective tasks observed differences with controls in regions implicated in salience, habit, goal-directed behavior, self-referential thinking and cognitive control. For nonaffective tasks, hyperactivity was observed in the insula, ACC and head of the caudate/putamen, while hypoactivity was observed in the medial prefrontal cortex (mPFC) and posterior caudate. Affective tasks were observed to relate to increased activation in the precuneus and posterior cingulate cortex, while decreased activation was found in the pallidum, ventral anterior thalamus and posterior caudate. The involvement of the cortico-striato-thalamo-cortical loop in OCD, as well as the high rates of comorbidity between OCD and ADHD, have led some to draw a link in their mechanism. Observed similarities include dysfunction of the anterior cingulate cortex and prefrontal cortex, as well as shared deficits in executive functions. The involvement of the orbitofrontal cortex and dorsolateral prefrontal cortex in OCD is shared with bipolar disorder and may explain the high degree of comorbidity. Decreased volumes of the dorsolateral prefrontal cortex related to executive function has also been observed in OCD.
People with OCD evince increased grey matter volumes in bilateral lenticular nuclei, extending to the caudate nuclei, with decreased grey matter volumes in bilateral dorsal medial frontal/anterior cingulate gyri. These findings contrast with those in people with other anxiety disorders, who evince decreased (rather than increased) grey matter volumes in bilateral lenticular/caudate nuclei, as well as decreased grey matter volumes in bilateral dorsal medial frontal/anterior cingulate gyri. Increased white matter volume and decreased fractional anisotropy in anterior midline tracts has been observed in OCD, possibly indicating increased fiber crossings.
Cognitive models
Generally, two categories of models for OCD have been postulated. The first category involves deficits in executive dysfunction and is based on the observed structural and functional abnormalities in the dlPFC, striatum and thalamus. The second category involves dysfunctional modulatory control and primarily relies on observed functional and structural differences in the ACC, mPFC and OFC.
One proposed model suggests that dysfunction in the orbitalfrontal cortex (OFC) leads to improper valuation of behaviors and decreased behavioral control, while the observed alterations in amygdala activations leads to exaggerated fears and representations of negative stimuli.
Due to the heterogeneity of OCD symptoms, studies differentiating various symptoms have been performed. Symptom-specific neuroimaging abnormalities include the hyperactivity of caudate and ACC in checking rituals, while finding increased activity of cortical and cerebellar regions in contamination-related symptoms. Neuroimaging differentiating content of intrusive thoughts has found differences between aggressive as opposed to taboo thoughts, finding increased connectivity of the amygdala, ventral striatum and ventromedial prefrontal cortex in aggressive symptoms, while observing increased connectivity between the ventral striatum and insula in sexual or religious intrusive thoughts.
Another model proposes that affective dysregulation links excessive reliance on habit-based action selection with compulsions. This is supported by the observation that those with OCD demonstrate decreased activation of the ventral striatum when anticipating monetary reward, as well as increased functional connectivity between the VS and the OFC. Furthermore, those with OCD demonstrate reduced performance in Pavlovian fear-extinction tasks, hyperresponsiveness in the amygdala to fearful stimuli and hyporesponsiveness in the amygdala when exposed to positively valanced stimuli. Stimulation of the nucleus accumbens has also been observed to effectively alleviate both obsessions and compulsions, supporting the role of affective dysregulation in generating both.
Neurobiological
From the observation of the efficacy of antidepressants in OCD, a serotonin hypothesis of OCD has been formulated. Studies of peripheral markers of serotonin, as well as challenges with proserotonergic compounds have yielded inconsistent results, including evidence pointing towards basal hyperactivity of serotonergic systems. Serotonin receptor and transporter binding studies have yielded conflicting results, including higher and lower serotonin receptor 5-HT2A and serotonin transporter binding potentials that were normalized by treatment with SSRIs. Despite inconsistencies in the types of abnormalities found, evidence points towards dysfunction of serotonergic systems in OCD. Orbitofrontal cortex overactivity is attenuated in people who have successfully responded to SSRI medication, a result believed to be caused by increased stimulation of serotonin receptors 5-HT2A and 5-HT2C.
A complex relationship between dopamine and OCD has been observed. Although antipsychotics, which act by antagonizing dopamine receptors, may improve some cases of OCD, they frequently exacerbate others. Antipsychotics, in the low doses used to treat OCD, may actually increase the release of dopamine in the prefrontal cortex, through inhibiting autoreceptors. Further complicating things is the efficacy of amphetamines, decreased dopamine transporter activity observed in OCD, and low levels of D2 binding in the striatum. Furthermore, increased dopamine release in the nucleus accumbens after deep brain stimulation correlates with improvement in symptoms, pointing to reduced dopamine release in the striatum playing a role in generating symptoms.
Abnormalities in glutamatergic neurotransmission have been implicated in OCD. Findings such as increased cerebrospinal glutamate, less consistent abnormalities observed in neuroimaging studies, and the efficacy of some glutamatergic drugs (such as the glutamate-inhibiting riluzole) have implicated glutamate in OCD. OCD has been associated with reduced N-Acetylaspartic acid in the mPFC, which is thought to reflect neuron density or functionality, although the exact interpretation has not been established.
Diagnosis
Formal diagnosis may be performed by a psychologist, psychiatrist, clinical social worker or other licensed mental health professional. OCD, like other mental and behavioral health disorders, cannot be diagnosed by a medical exam, nor are there any medical exams that can predict if one will fall victim to such illnesses. To be diagnosed with OCD, a person must have obsessions, compulsions or both, according to the Diagnostic and Statistical Manual of Mental Disorders (DSM). The DSM notes that there are multiple characteristics that can turn obsessions and compulsions from normalized behavior to "clinically significant". There has to be recurring and strong thoughts or impulsive that intrude on the day-to-day lives of the patients and cause noticeable levels of anxiousness.
These thoughts, impulses or images are of a degree or type that lies outside the normal range of worries about conventional problems. A person may attempt to ignore or suppress such obsessions, neutralize them with another thought or action, or try to rationalize their anxiety away. People with OCD tend to recognize their obsessions as irrational.
Compulsions become clinically significant when a person feels driven to perform them in response to an obsession or according to rules that must be applied rigidly and when the person consequently feels or causes significant distress. Therefore, while many people who do not have OCD may perform actions often associated with OCD (such as ordering items in a pantry by height), the distinction with clinically significant OCD lies in the fact that the person with OCD must perform these actions to avoid significant psychological distress. These behaviors or mental acts are aimed at preventing or reducing distress or preventing some dreaded event or situation; however, these activities are not logically or practically connected to the issue, or, they are excessive.
Moreover, the obsessions or compulsions must be time-consuming, often taking up more than one hour per day or cause impairment in social, occupational or scholastic functioning. It is helpful to quantify the severity of symptoms and impairment before and during treatment for OCD. In addition to the person's estimate of the time spent each day harboring obsessive-compulsive thoughts or behaviors, concrete tools can be used to gauge the person's condition. This may be done with rating scales, such as the Yale–Brown Obsessive Compulsive Scale (Y-BOCS; expert rating) or the obsessive–compulsive inventory (OCI-R; self-rating). With measurements such as these, psychiatric consultation can be more appropriately determined, as it has been standardized.
In regards to diagnosing, the health professional also looks to make sure that the signs of obsessions and compulsions are not the results of any drugs, prescription or recreational, that the patient may be taking.
There are several types of obsessive thoughts that are found commonly in those with OCD. Some of these include fear of germs, hurting loved ones, embarrassment, neatness, societally unacceptable sexual thoughts etc. Within OCD, these specific categories are often diagnosed into their own type of OCD.
OCD is sometimes placed in a group of disorders called the obsessive–compulsive spectrum.
Another criterion in the DSM is that a person's mental illness does not fit one of the other categories of a mental disorder better. That is to say, if the obsessions and compulsions of a patient could be better described by trichotillomania, it would not be diagnosed as OCD. That being said, OCD does often go hand in hand with other mental disorders. For this reason, one may be diagnosed with multiple mental disorders at once.
A different aspect of the diagnoses is the degree of insight had by the individual in regards to the truth of the obsessions. There are three levels, good/fair, poor and absent/delusional. Good/fair indicated that the patient is aware that the obsessions they have are not true or probably not true. Poor indicates that the patient believes their obsessional beliefs are probably true. Absent/delusional indicates that they fully believe their obsessional thoughts to be true. Approximately 4% or fewer individuals with OCD will be diagnosed as absent/delusional. Additionally, as many as 30% of those with OCD also have a lifetime tic disorder, meaning they have been diagnosed with a tic disorder at some point in their life.
There are several different types of tics that have been observed in individuals with OCD. These include but are not limited to, "grunting", "jerking" or "shrugging" body parts, sniffling and excessive blinking.
There has been a significant amount of progress over the last few decades and as of 2022 there is statically significant improvement in the diagnostic process for individuals with OCD. One study found that of two groups of individuals, one with participants under the age of 27.25 and one with participants over that age, those in the younger group experienced a significantly faster time between the onset of OCD tendencies and their formal diagnoses.
Differential diagnosis
OCD is often confused with the separate condition obsessive–compulsive personality disorder (OCPD). OCD is egodystonic, meaning that the disorder is incompatible with the individual's self-concept. As egodystonic disorders go against a person's self-concept, they tend to cause much distress. OCPD, on the other hand, is egosyntonic, marked by the person's acceptance that the characteristics and behaviors displayed as a result are compatible with their self-image, or are otherwise appropriate, correct or reasonable.
As a result, people with OCD are often aware that their behavior is not rational and are unhappy about their obsessions, but nevertheless feel compelled by them. By contrast, people with OCPD are not aware of anything abnormal; they will readily explain why their actions are rational. It is usually impossible to convince them otherwise and they tend to derive pleasure from their obsessions or compulsions.
Management
Cognitive behavioral therapy (CBT) and psychotropic medications are the first-line treatments for OCD.
Therapy
One specific CBT technique used is called exposure and response prevention (ERP), which involves teaching the person to deliberately come into contact with situations that trigger obsessive thoughts and fears (exposure), without carrying out the usual compulsive acts associated with the obsession (response prevention). This technique causes patients to gradually learn to tolerate the discomfort and anxiety associated with not performing their compulsions. For many patients, ERP is the add-on treatment of choice when selective serotonin reuptake inhibitors (SSRIs) or serotonin–norepinephrine reuptake inhibitors (SNRIs) medication does not effectively treat OCD symptoms, or vice versa, for individuals who begin treatment with psychotherapy. This technique is considered superior to others due to the lack of medication used. However, up to 25% of patients will discontinue treatment due to the severity of their tics. CBT normally lasts anywhere from 12-16 sessions, with homework assigned to the patient in between meetings with a therapist. (Lack 2012). Modalities differ in ERP treatment but both virtual reality based as well as unguided computer assisted treatment programs have shown effective results in treatment programs.
For example, a patient might be asked to touch something very mildly contaminated (exposure) and wash their hands only once afterward (response prevention). Another example might entail asking the patient to leave the house and check the lock only once (exposure), without going back to check again (response prevention). After succeeding at one stage of treatment, the patient's level of discomfort in the exposure phase can be increased. When this therapy is successful, the patient will quickly habituate to an anxiety-producing situation, discovering a considerable drop in anxiety level.
ERP has a strong evidence base and is considered the most effective treatment for OCD. However, this claim was doubted by some researchers in 2000, who criticized the quality of many studies. While ERP can lead a majority of clients to improvements, many do not reach remission or become asymptomatic; some therapists are also hesitant to use this approach.
The recent development of remotely technology-delivered CBT is increasing access to therapy options for those living with OCD and remote versions appear to equally as effective as in-person therapy options. The development of smartphone interventions for OCD that utilize CBT techniques are another alternative that is expanding access to therapy while allowing therapies to be personalized for each patient.
Acceptance and commitment therapy (ACT), a newer therapy also used to treat anxiety and depression, has also been found to be effective in treatment of OCD. ACT uses acceptance and mindfulness strategies to teach patients not to overreact to or avoid unpleasant thoughts and feelings but rather "move toward valued behavior".
Inference-based therapy (IBT) is a form of cognitive therapy specifically developed for treating OCD. The therapy posits that individuals with OCD put a greater emphasis on an imagined possibility than on what can be perceived with the senses, and confuse the imagined possibility with reality, in a process called inferential confusion. According to inference-based therapy, obsessional thinking occurs when the person replaces reality and real probabilities with imagined possibilities. The goal of inference-based therapy is to reorient clients towards trusting the senses and relating to reality in a normal, non-effortful way. Differences between normal and obsessional doubts are presented and clients are encouraged to use their senses and reasoning as they do in non-obsessive–compulsive disorder situations. Research on Inference-Based Cognitive-Behavior Therapy (I-CBT) suggests it can lead to improvements for those with OCD.
A 2007 Cochrane review found that psychological interventions derived from CBT models, such as ERP, ACT and IBT, were more effective than non-CBT interventions. Other forms of psychotherapy, such as psychodynamics and psychoanalysis, may help in managing some aspects of the disorder. However, in 2007, the American Psychiatric Association (APA) noted a lack of controlled studies showing their efficacy, "in dealing with the core symptoms of OCD". For body-focused repetitive behaviors (BFRB), behavioral interventions such as habit-reversal training and decoupling are recommended.
Psychotherapy in combination with psychiatric medication may be more effective than either option alone for individuals with severe OCD. ERP coupled with weight restoration and serotonin reuptake inhibitors has proven the most effective when treating OCD and an eating disorder simultaneously.
Medication
The medications most frequently used to treat OCD are antidepressants, including selective serotonin reuptake inhibitors (SSRIs) and serotonin–norepinephrine reuptake inhibitors (SNRIs). Sertraline and fluoxetine are effective in treating OCD for children and adolescents.
SSRIs are a second-line treatment of adult OCD with mild functional impairment and as first-line treatment for those with moderate or severe impairment. In children, SSRIs can be considered as a second-line therapy in those with moderate to severe impairment, with close monitoring for psychiatric adverse effects. Patients treated with SSRIs are about twice as likely to respond to treatment as are those treated with placebo, so this treatment is qualified as efficacious. Efficacy has been demonstrated both in short-term (6–24 weeks) treatment trials and in discontinuation trials with durations of 28–52 weeks.
Clomipramine, a medication belonging to the class of tricyclic antidepressants, appears to work as well as SSRIs, but has a higher rate of side effects.
In 2006, the National Institute for Health and Care Excellence (NICE) guidelines recommended augmentative second-generation (atypical) antipsychotics for treatment-resistant OCD. Atypical antipsychotics are not useful when used alone and no evidence supports the use of first-generation antipsychotics. For OCD treatment specifically, there is tentative evidence for risperidone and insufficient evidence for olanzapine. Quetiapine is no better than placebo with regard to primary outcomes, but small effects were found in terms of Y-BOCS score. The efficacy of quetiapine and olanzapine are limited by an insufficient number of studies. A 2014 review article found two studies that indicated that aripiprazole was "effective in the short-term" and found that "[t]here was a small effect-size for risperidone or antipsychotics in general in the short-term"; however, the study authors found "no evidence for the effectiveness of quetiapine or olanzapine in comparison to placebo." While quetiapine may be useful when used in addition to an SSRI/SNRI in treatment-resistant OCD, these drugs are often poorly tolerated and have metabolic side effects that limit their use. A guideline by the American Psychological Association suggested that dextroamphetamine may be considered by itself after more well-supported treatments have been attempted.
Procedures
Electroconvulsive therapy (ECT) has been found to have effectiveness in some severe and refractory cases. Transcranial magnetic stimulation has shown to provide therapeutic benefits in alleviating symptoms.
Surgery may be used as a last resort in people who do not improve with other treatments. In this procedure, a surgical lesion is made in an area of the brain (the cingulate cortex). In one study, 30% of participants benefitted significantly from this procedure. Deep brain stimulation and vagus nerve stimulation are possible surgical options that do not require destruction of brain tissue. However, because deep brain stimulation results in such an instant and intense change, individuals may experience identity challenges afterward. In the United States, the Food and Drug Administration (FDA) approved deep brain stimulation for the treatment of OCD under a humanitarian device exemption, requiring that the procedure be performed only in a hospital with special qualifications to do so.
In the United States, psychosurgery for OCD is a treatment of last resort and will not be performed until the person has failed several attempts at medication (at the full dosage) with augmentation, and many months of intensive cognitive behavioral therapy with exposure and ritual/response prevention. Likewise, in the United Kingdom, psychosurgery cannot be performed unless a course of treatment from a suitably qualified cognitive–behavioral therapist has been carried out.
Children
Therapeutic treatment may be effective in reducing ritual behaviors of OCD for children and adolescents. Similar to the treatment of adults with OCD, cognitive behavioral therapy, along with exposure and response prevention (ERP) therapy, stands as an effective and validated first line of treatment of OCD in children. Family involvement, in the form of behavioral observations and reports, is a key component to the success of such treatments. Parental interventions also provide positive reinforcement for a child who exhibits appropriate behaviors as alternatives to compulsive responses. In a recent meta-analysis of evidenced-based treatment of OCD in children, family-focused individual CBT was labeled as "probably efficacious", establishing it as one of the leading psychosocial treatments for youth with OCD. After one or two years of therapy, in which a child learns the nature of their obsession and acquires strategies for coping, they may acquire a larger circle of friends, exhibit less shyness and become less self-critical. Trials have shown that children and adolescents with OCD should begin treatment with the combination of CBT with a selective serotonin reuptake inhibitor or CBT alone, rather than only an SSRI. A 2024 systeramitc review of the literature found that combining ERP therapy with selective serotonin reuptake inhibitors can enhance treatment outcomes compared to using SSRIs alone.
Although the known causes of OCD in younger age groups range from brain abnormalities to psychological preoccupations, life stress such as bullying and traumatic familial deaths may also contribute to childhood cases of OCD, and acknowledging these stressors can play a role in treating the disorder.
Prognosis
Quality of life is reduced across all domains in OCD. While psychological or pharmacological treatment can lead to a reduction of OCD symptoms and an increase in reported quality of life, symptoms may persist at moderate levels even following adequate treatment courses, and completely symptom-free periods are uncommon. In pediatric OCD, around 40% still have the disorder in adulthood and around 40% qualify for remission. The risk of having at least one comorbid personality disorder in OCD is 52%, which is the highest among anxiety disorders and greatly impacts its management and prognosis.
Epidemiology
Obsessive–compulsive disorder affects about 2.3% of people at some point in their life, with the yearly rate about 1.2%. OCD occurs worldwide. It is unusual for symptoms to begin after the age of 35 and half of people develop problems before 20. Males and females are affected about equally. However, there is an earlier age for onset for males than females.
History
Plutarch, an ancient Greek philosopher and historian, describes an ancient Roman man who possibly had scrupulosity, which could be a symptom of OCD or OCPD. This man is described as "turning pale under his crown of flowers", praying with a "faltering voice" and scattering "incense with trembling hands".
In the 7th century AD, John Climacus records an instance of a young monk plagued by constant and overwhelming "temptations to blasphemy" consulting an older monk, who told him: "My son, I take upon myself all the sins which these temptations have led you, or may lead you, to commit. All I require of you is that for the future you pay no attention to them whatsoever." The Cloud of Unknowing, a Christian mystical text from the late 14th century, recommends dealing with recurring obsessions by attempting to ignore them, and, if that fails, to "cower under them like a poor wretch and a coward overcome in battle, and reckon it to be a waste of your time for you to strive any longer against them", a technique now known as emotional flooding.
Abu Zayd Al-Balkhi, the 9th century Islamic polymath, was likely the first to classify OCD into different types and pioneer cognitive behavioral therapy, in a fashion unique to his era and which was not popular in Greek medicine. In his medical treatise entitled Sustenance of the Body and Soul, Al-Balkhi describes obsessions particular to the disorder as "Annoying thoughts that are not real. These intrusive thoughts prevent enjoying life, and performing daily activities. They affect concentration and interfere with ability to carry out different tasks." As treatment, Al-Balkhi suggests treating obsessive thoughts with positive thoughts and mind-based therapy.
From the 14th to the 16th century in Europe, it was believed that people who experienced blasphemous, sexual or other obsessive thoughts were possessed by the devil. Based on this reasoning, treatment involved banishing the "evil" from the "possessed" person through exorcism. The vast majority of people who thought that they were possessed by the devil did not have hallucinations or other "spectacular symptoms" but "complained of anxiety, religious fears, and evil thoughts." In 1584, a woman from Kent, England, named Mrs. Davie, described by a justice of the peace as "a good wife", was nearly burned at the stake after she confessed that she experienced constant, unwanted urges to murder her family.
The English term obsessive–compulsive arose as a translation of German Zwangsvorstellung (obsession) used in the first conceptions of OCD by Karl Westphal. Westphal's description went on to influence Pierre Janet, who further documented features of OCD. In the early 1910s, Sigmund Freud attributed obsessive–compulsive behavior to unconscious conflicts that manifest as symptoms. Freud describes the clinical history of a typical case of "touching phobia" as starting in early childhood, when the person has a strong desire to touch an item. In response, the person develops an "external prohibition" against this type of touching. However, this "prohibition does not succeed in abolishing" the desire to touch; all it can do is repress the desire and "force it into the unconscious." Freudian psychoanalysis remained the dominant treatment for OCD until the mid-1980s, even though medicinal and therapeutic treatments were known and available, because it was widely thought that these treatments would be detrimental to the effectiveness of the psychotherapy. In the mid-1980s, this approach changed and practitioners began treating OCD primarily with medicine and practical therapy rather than through psychoanalysis.
One of the first successful treatments of OCD, exposure and response prevention, emerged during the 1960s, when psychologist Vic Meyer exposed two hospitalized patients to anxiety-inducing situations while preventing them from performing any compulsions. Eventually, both patients' anxiety level dropped to manageable levels. Meyer devised this procedure from his analysis of fear extinguishment in animals via flooding. The success of ERP clinically and scientifically has been summarized as "spectacular" by prominent OCD researcher Stanley Rachman decades following Meyer's creation of the method.
In 1967, psychiatrist Juan José López-Ibor reported that the drug clomipramine was effective in treating OCD. Many reports of its success in treatment followed and several studies had confirmed its effectiveness by the 1980s. However, clomipramine was subsequently displaced by new SSRIs developed in the 1970s, such as fluoxetine and sertraline, which were shown to have fewer side effects.
Obsessive–compulsive symptoms worsened during the early stages of the COVID-19 pandemic, particularly for individuals with contamination-related OCD.
Notable cases
John Bunyan (1628–1688), the author of The Pilgrim's Progress, displayed symptoms of OCD (which had not yet been named). During the most severe period of his condition, he would mutter the same phrase over and over again to himself while rocking back and forth. He later described his obsessions in his autobiography Grace Abounding to the Chief of Sinners, stating, "These things may seem ridiculous to others, even as ridiculous as they were in themselves, but to me they were the most tormenting cogitations." He wrote two pamphlets advising those with similar anxieties. In one of them, he warns against indulging in compulsions: "Have care of putting off your trouble of spirit in the wrong way: by promising to reform yourself and lead a new life, by your performances or duties."
British poet, essayist and lexicographer Samuel Johnson (1709–1784) also had OCD. He had elaborate rituals for crossing the thresholds of doorways and repeatedly walked up and down staircases counting the steps. He would touch every post on the street as he walked past, only step in the middle of paving stones and repeatedly perform tasks as though they had not been done properly the first time.
The "Rat Man", real name Ernst Lanzer, a patient of Sigmund Freud, suffered from what was then called "obsessional neurosis". Lanzer's illness was characterised most famously by a pattern of distressing intrusive thoughts in which he feared that his father or a female friend would be subjected to a purported Chinese method of torture in which rats would be encouraged to gnaw their way out of a victim's body by a hot poker.
American aviator and filmmaker Howard Hughes is known to have had OCD, primarily an obsessive fear of germs and contamination. Friends of Hughes have also mentioned his obsession with minor flaws in clothing. This was conveyed in The Aviator (2004), a film biography of Hughes.
English singer-songwriter George Ezra has openly spoken about his life-long struggle with OCD, particularly primarily obsessional obsessive–compulsive disorder (Pure O).
Swedish climate activist Greta Thunberg is also known to have OCD, among other mental health conditions.
American actor James Spader has also spoken about his OCD. In 2014, when interviewed for Rolling Stone he said: "I'm obsessive-compulsive. I have very, very strong obsessive-compulsive issues. I'm very particular. ... It's very hard for me, you know? It makes you very addictive in behavior, because routine and ritual become entrenched. But in work, it manifests itself in obsessive attention to detail and fixation. It serves my work very well: Things don't slip by. But I'm not very easygoing.
In 2022 the president of Chile Gabriel Boric stated that he had OCD, saying: "I have an obsessive–compulsive disorder that's completely under control. Thank God I've been able to undergo treatment and it doesn't make me unable to carry out my responsibilities as the President of the Republic."
In a documentary released in 2023, David Beckham shared details about his compelling cleaning rituals, need for symmetry in the fridge and the impact of OCD on his life.
Society and culture
Art, entertainment and media
Movies and television shows may portray idealized or incomplete representations of disorders such as OCD. Compassionate and accurate literary and on-screen depictions may help counteract the potential stigma associated with an OCD diagnosis and lead to increased public awareness, understanding and sympathy for such disorders.
The play and film adaptations of The Odd Couple based around the character of Felix, who shows some of the common symptoms of OCD.
In the film As Good as It Gets (1997), actor Jack Nicholson portrays a man with OCD who performs ritualistic behaviors that disrupt his life.
The film Matchstick Men (2003) portrays a con man named Roy (Nicolas Cage) with OCD who opens and closes doors three times while counting aloud before he can walk through them.
In the television series Monk (2002–2009), the titular character Adrian Monk fears both human contact and dirt.
The one-man show The Life and Slimes of Marc Summers (2016), a stage adaptation of Marc Summers' 1999 memoir which recounts how OCD affected his entertainment career.
In the novel Turtles All the Way Down (2017) by John Green, teenage main character Aza Holmes struggles with OCD that manifests as a fear of the human microbiome. Throughout the story, Aza repeatedly opens an unhealed callus on her finger to drain out what she believes are pathogens. The novel is based on Green's own experiences with OCD. He explained that Turtles All the Way Down is intended to show how "most people with chronic mental illnesses also live long, fulfilling lives."
The British TV series Pure (2019) stars Charly Clive as a 24-year-old Marnie who is plagued by disturbing sexual thoughts, as a kind of primarily obsessional obsessive compulsive disorder.
Research
The naturally occurring sugar inositol has been suggested as a treatment for OCD.
μ-Opioid receptor agonists, such as hydrocodone and tramadol, may improve OCD symptoms. Administration of opioids may be contraindicated in individuals concurrently taking CYP2D6 inhibitors such as fluoxetine and paroxetine.
Much research is devoted to the therapeutic potential of the agents that affect the release of the neurotransmitter glutamate or the binding to its receptors. These include riluzole, memantine, gabapentin, N-acetylcysteine (NAC), topiramate and lamotrigine. Research on the potential for other supplements, such as milk thistle, to help with OCD and various neurological disorders, is ongoing.
Researchers have identified over 600 genes related to cortical thickness, a factor that impacts OCD expression. "Notably, the enrichment of genes involved in ion transport regulation, responses to environmental stimuli, and metal ion transport regulation suggests the roles of these processes in OCD pathophysiology."
Research indicates that people with OCD have a lower amplitude of low-frequency fluctuation in both the left and right putamen. The right putamen also displays decreased functional connectivity with the left putamen which extends to the left inferior frontal gyrus (IFG), bilateral precuneus extending to calcarine, right middle occipital cortex extending to the right middle temporal cortex, and left middle occipital gyrus. In addition, the decreased connectivity between the right putamen and the left putamen is negatively correlated with Y-BOCS scores.
In a study exploring the correlation between neural biomarkers and response to transcranial Direct Current Stimulation (tDCS) in people with OCD, researchers found thicker precentral and paracentral areas in people with OCD compared to controls. A significant association was found between a thinner precentral area and reduced YBOCS scores.
Other animals
Advocacy
Many organizations and charities around the world advocate for the wellbeing of people with OCD, stigma reduction, research and awareness. The International OCD Foundation (IOCDF) is the largest 501(c)3 nonprofit organization dedicated to serving a broad community of individuals with OCD and related disorders, their family members and loved ones, and mental health professionals and researchers around the world. Since 1986, the IOCDF provides up-to-date education and resources, strengthens community engagement worldwide, delivers quality professional training to clinicians and funds groundbreaking research.
| Biology and health sciences | Mental disorders | Health |
6587493 | https://en.wikipedia.org/wiki/Candareen | Candareen | A candareen (; ; Singapore English usage: hoon) is a traditional measurement of weight in East Asia. It is equal to 10 cash and is of a mace. It is approximately 378 milligrams. A troy candareen is approximately .
In Hong Kong, one candareen is 0.3779936375 grams and, in the Weights and Measures Ordinance, it is ounces avoirdupois. In Singapore, one candareen is 0.377994 grams.
The word candareen comes from the Malay kandūri. An earlier English form of the name was condrin. The candareen was also formerly used to describe a unit of currency in imperial China equal to 10 li () and is of a mace. The Mandarin Chinese word fēn is used to denote of a Chinese renminbi yuan but the term candareen for that currency is now obsolete.
Postal denomination
On 1 May 1878 the Imperial Maritime Customs was opened to the public and China's first postage stamps, the "Large Dragons" (), were issued to handle payment. The stamps were inscribed "CHINA" in both Latin and Chinese characters, and denominated in candareens.
| Physical sciences | Chinese | Basics and measurement |
6587584 | https://en.wikipedia.org/wiki/Mace%20%28unit%29 | Mace (unit) | A mace (; Hong Kong English usage: tsin; Southeast Asian English usage: chee) is a traditional Chinese measurement of weight in East Asia that was also used as a currency denomination. It is equal to 10 candareens and is of a tael or approximately 3.78 grams. A troy mace is approximately 3.7429 grams. In Hong Kong, one mace is grams. and in Ordinance 22 of 1884, it is ounces avoirdupois. In Singapore, one mace (referred to as chee) is grams.
In imperial China, 10 candareens equaled 1 mace which was of a tael and, like the other units, was used in weight-denominated silver currency system. A common denomination was 7 mace and 2 candareens, equal to one silver Chinese yuan.
Name
Like other similar measures such as tael and catty, the English word "mace" derives from Malay, in this case through Dutch maes, plural masen, from Malay mas which, in turn, derived from Sanskrit (), a word related to "mash," another name for the urad bean, and masha, a traditional Indian unit of weight equal to 0.97 gram. This word is unrelated to other uses of "mace" in English.
The Chinese word for mace is qián (), which is also a generic word for "money" in Mandarin Chinese. The same Chinese character (kanji) was used for the Japanese sen, the former unit equal to of a Japanese yen, the Korean chŏn (revised: jeon), the former unit equal to of a Korean won, and for the Vietnamese tiền, a currency used in late imperial Vietnam, although none of these has ever been known as "mace" in English.
| Physical sciences | Chinese | Basics and measurement |
6588022 | https://en.wikipedia.org/wiki/Lithium%20borohydride | Lithium borohydride | Lithium borohydride (LiBH4) is a borohydride and known in organic synthesis as a reducing agent for esters. Although less common than the related sodium borohydride, the lithium salt offers some advantages, being a stronger reducing agent and highly soluble in ethers, whilst remaining safer to handle than lithium aluminium hydride.
Preparation
Lithium borohydride may be prepared by the metathesis reaction, which occurs upon ball-milling the more commonly available sodium borohydride and lithium bromide:
NaBH4 + LiBr → NaBr + LiBH4
Alternatively, it may be synthesized by treating boron trifluoride with lithium hydride in diethyl ether:
BF3 + 4 LiH → LiBH4 + 3 LiF
Reactions
Lithium borohydride is useful as a source of hydride (H–). It can react with a range of carbonyl substrates and other polarized carbon structures to form a hydrogen–carbon bond. It can also react with Brønsted–Lowry-acidic substances (sources of H+) to form hydrogen gas.
Reduction reactions
As a hydride reducing agent, lithium borohydride is stronger than sodium borohydride but weaker than lithium aluminium hydride. Unlike the sodium analog, it can reduce esters to alcohols, nitriles and primary amides to amines, and can open epoxides. The enhanced reactivity in many of these cases is attributed to the polarization of the carbonyl substrate by complexation to the lithium cation. Unlike the aluminium analog, it does not react with nitro groups, carbamic acids, alkyl halides, or secondary and tertiary amides.
Hydrogen generation
Lithium borohydride reacts with water to produce hydrogen. This reaction can be used for hydrogen generation.
Although this reaction is usually spontaneous and violent, somewhat-stable aqueous solutions of lithium borohydride can be prepared at low temperature if degassed, distilled water is used and exposure to oxygen is carefully avoided.
Energy storage
Lithium borohydride is renowned as one of the highest-energy-density chemical energy carriers. Although presently of no practical importance, the solid liberates 65 MJ/kg heat upon treatment with atmospheric oxygen. Since it has a density of 0.67 g/cm3, oxidation of liquid lithium borohydride gives 43 MJ/L. In comparison, gasoline gives 44 MJ/kg (or 35 MJ/L), while liquid hydrogen gives 120 MJ/kg (or 8.0 MJ/L). The high specific energy density of lithium borohydride has made it an attractive candidate to propose for automobile and rocket fuel, but despite the research and advocacy, it has not been used widely. As with all chemical-hydride-based energy carriers, lithium borohydride is very complex to recycle (i.e. recharge) and therefore suffers from a low energy conversion efficiency. While batteries such as lithium-ion carry an energy density of up to 0.72 MJ/kg and 2.0 MJ/L, their DC-to-DC conversion efficiency can be as high as 90%. In view of the complexity of recycling mechanisms for metal hydrides, such high energy-conversion efficiencies are not practical with present technology.
| Physical sciences | Borohydride salts | Chemistry |
1340551 | https://en.wikipedia.org/wiki/Crown%20%28dental%20restoration%29 | Crown (dental restoration) | In dentistry, a crown or a dental cap is a type of dental restoration that completely caps or encircles a tooth or dental implant. A crown may be needed when a large dental cavity threatens the health of a tooth. Some dentists will also finish root canal treatment by covering the exposed tooth with a crown. A crown is typically bonded to the tooth by dental cement. They can be made from various materials, which are usually fabricated using indirect methods. Crowns are used to improve the strength or appearance of teeth and to halt deterioration. While beneficial to dental health, the procedure and materials can be costly.
The most common method of crowning a tooth involves taking a dental impression of a tooth prepared by a dentist, then fabricating the crown outside of the mouth. The crown can then be inserted at a subsequent dental appointment. This indirect method of tooth restoration allows use of strong restorative material requiring time-consuming fabrication under intense heat, such as casting metal or firing porcelain, that would not be possible inside the mouth. Because of its compatible thermal expansion, relatively similar cost, and cosmetic difference, some patients choose to have their crown fabricated with gold.
Computer technology is increasingly employed for crown fabrication in CAD/CAM dentistry.
Indications for dental crowns
Crowns are indicated to:
Replace existing crowns which have failed.
Restore the form, function and appearance of badly broken down, worn or fractured teeth, where other simpler forms of restorations are unsuitable or have been found to fail clinically.
Improve the aesthetics of unsightly teeth which cannot be managed by simpler cosmetic and restorative procedures.
Maintain the structural stability and reduce the risk of fractures of extensively restored teeth including those which have been endodontically treated.
Restore the visible portion of a single dental implant.
Restoration of endodontically treated teeth
Traditionally, it has been proposed that teeth which have undergone root canal treatment are more likely to fracture and therefore require cuspal protection by providing occlusal coverage with an indirect restoration like crowns. This led to routine prescribing of crowns for root-treated teeth. However, recent review of literature reveals that there is no strong evidence to show that crowns are better than other routine restorations to restore root-filled teeth. The general advice is that dentists should use their clinical experience in view of the patient's preferences when making the decision of using a crown. As a rule of thumb, the use of crowns and other indirect restorations for root treated teeth is justified when the surface area of the access cavity exceeds one third of the occlusal surface of the tooth, when the lingual or buccal walls are undermined or when the mesial and distal marginal ridges are missing.
Clinical stages of dental crown provision
Assessment
Choice of restoration
Tooth preparation
Construction and fit of temporary restoration
Tooth preparation impressions
Fit of definitive restoration
Short-term follow up
Long-term follow up
Assessment
In order to ensure optimum condition and longevity for the proposed crowns, several factors need to be explored by conducting a thorough and targeted patient history and clinical dental examination. These factors include:
Patient factors
Patient expectations
Patient motivation to adhere to the treatment plan and maintain results
Financial and time costs to the patient
Biological factors
Periodontal health status and periodontal disease risk
Pulpal health and endodontic disease risk
Caries and caries risk
Occlusion and occlusal problems risk
Mechanical factors
Amount of remaining tooth structure
Height and width of tooth to be prepared
Attachment levels of the tooth to be prepared
Root shape and length of the tooth to be prepared
Aesthetic factors
Choice of restoration
The choice(s) of crown restoration can be described by:
The dimensions and percentage coverage of the natural crown
Full crowns
3/4 and 7/8 crowns
Material to be used
Metal
Metal-ceramic crowns
Full ceramic crowns
3/4 and 7/8 crowns
These restorations are a hybrid between an onlay and a full crown. They are named based on the estimated wall coverage of the walls of the tooth; e.g. the 3/4 crown aims to cover three out of the four walls, with the buccal wall being usually spared, thus reducing sound tooth tissue to be prepared. They are normally fabricated in gold. Grooves or boxes are normally added to the preparations as close to the unprepared wall as possible to increase retention of the crown. Despite its advantages of reducing sound tooth preparation, these crowns are not commonly prescribed in practice because they are technically difficult and have poor patient acceptability due to the metal showing through in their smile.
Full metal crowns
As the name suggests, these crowns are entirely cast in a metal alloy. There are a multitude of alloys available and the selection of a particular alloy over another depends on several factors including cost, handling, physical properties, biocompatibility. The American Dental Association categories alloys in three groups: high-noble, noble and base metal alloys.
High-noble and noble alloys
Noble and high-noble alloys used in casting crowns are generally based on alloys of gold. Gold is not used in its pure form as it is too soft and has poor mechanical strength. Other metals included in gold alloys are copper, platinum, palladium, zinc, indium and nickel. All types of gold casting alloys used in prosthodontics (Type I - IV) are categorised by their percentage content of gold and hardness, with Type I being the softest and Type IV the hardest. Generally, Type III and IV alloys (62 - 78% and 60 - 70% gold content respectively) are used in casting of full crowns, as these are hard enough to withstand occlusal forces. Gold crowns (also known as gold shell crowns) are generally indicated for posterior teeth due to aesthetic reasons. They are durable in function and strong in thin sections, therefore require minimal tooth preparation. They also have similar wear properties to enamel, so they are not likely to cause excessive wear to the opposing tooth. They have good dimensional accuracy when cast which minimises chair-side/appointment time and can be relatively easy to polish if any changes are required. Palladium based alloys are also used. These were introduced as a cheaper alternative to gold alloys in the 1970s. Palladium has a strong whitening effect giving most of its alloys a silverish appearance.
Base-metal alloys
Cast base-metal alloys are rarely used to make full metal crowns. They are more commonly used as part of metal-ceramic crowns as bonding alloys. When compared to high-noble and noble alloys, they are stronger and harder; they can be used in thinner sections (0.3mm as opposed to 0.5mm) however they are harder to adjust and are more likely to cause excessive wear on real opposing teeth. Furthermore, there may be problems with people who have a nickel allergy.
Common base-metal alloys used in dentistry are:
Silver-palladium
Silver-palladium-copper
Nickel-chromium
Nickel-chromium-beryllium
Cobalt-chromium
Titanium
Titanium
Titanium and titanium alloys are highly biocompatible. Its strength, rigidity and ductility are similar to that of other casting alloys used in dentistry. Titanium also readily forms an oxide layer on its surface which gives it anti-corrosive properties and allows it to bond to ceramics, a useful property in the manufacture of metal-ceramic crowns.
Full ceramic crowns
Dental ceramics or porcelains are used for crowns manufacture primarily for their aesthetic properties compared to all metal restorations. These materials are generally quite brittle and prone to fracture. Many classifications have been used to categorise dental ceramics, with the simplest, based on the material from which they are made, i.e. silica, alumina or zirconia.
Silica
Silica-based ceramics are highly aesthetic due to their high glass content and excellent optical properties due to the addition of filler particles which enhance opalescence, fluorescence which can mimic the colour of natural enamel and dentine. These ceramics, however, suffer from poor mechanical strength, and therefore often used for veneering stronger substructures.
Examples include aluminosilicate glass, e.g. feldspathic, synthetic porcelain, and leucite reinforced ceramics.
Mechanical properties can be improved by the addition of filler particles, e.g. lithium disilicate, and are therefore termed glass ceramics. Glass-ceramics can be used alone to make all-ceramic restorations either as a single form (termed uni-layered) or can act as a substructures for subsequent veneering (or layering) with weaker feldspathic porcelain (restorations termed bi-layered).
Alumina
Alumina (aluminium oxide) was introduced as a dental substructure (core) in 1989 when the material was slip cast, sintered, and infiltrated with glass. More recently, glass-infiltrated alumina cores are produced by electrophoretic deposition, a rapid nanofabrication process. During this process, particles of a slip are brought to the surface of a dental die by an electric current, thereby forming a precision-fitting core greenbody in seconds. Margins are then trimmed and the greenbody is sintered and infiltrated with glass. Glass-infiltrated alumina has significantly higher porcelain bond strength over CAD/CAM produced zirconia and alumina cores without glass.
Alumina cores without glass are produced by milling pre-sintered blocks of the material utilizing a CAD/CAM dentistry technique. Cores without glass must be oversized to compensate for shrinkage that occurs when the core is fully sintered. Milled cores are then sintered and shrink to the correct size.
All alumina cores are layered with tooth tissue-like feldspathic porcelain to make true-to-life color and shape. Dental artists called ceramists, can customize the "look" of these crowns to individual patient and dentist requirements. Alumina cores have better translucency than zirconia, but worse than lithium disilicate.
Zirconia
Yttria-stabilized zirconia, also known simply as zirconia, is a very hard ceramic that is used as a strong base material in some full ceramic restorations. Zirconia is relatively new in dentistry and the published clinical data is correspondingly limited. The zirconia used in dentistry is zirconium oxide (ZrO2) which has been stabilized with the addition of yttrium oxide. Yttria-stabilized zirconia is also known as YSZ.
The zirconia substructure (core) is usually designed on a digital representation of the patient's mouth, which is captured with a three-dimensional digital scan of the patient, impression, or model. The core is then milled from a block of zirconia in a soft pre-sintered state. Once milled, the zirconia is sintered in a furnace where it shrinks by 20% and reaches its full strength of 850–1000 MPa. Recently, the strength of zirconia for dental restorations reaching 1200 MPa is reported. The zirconia core structure can be layered with tooth tissue-like feldspathic porcelain to create the final color and shape of the tooth. Because bond strength of layered porcelain fused to zirconia is not strong; chipping of the conventional veneering ceramic frequently occurs, crowns and bridges are nowadays increasingly made with monolithic zirconia crowns produced from a color and structure graded zirconia block, and coated with a thin layer of glaze stains. Esthetic prosthetic restorations, with natural reflection, color from within and color gradients influenced by the internal dentinal core anatomy can best be accomplished by veneered zirconia, rather than with crowns of monolithic zirconia. In the production of dental restorations specifically made for one patient, dental technicians with their problem-solving skills, dexterity and cognitive skills are until recently the only way to provide the required esthetics, individuality and artistry with porcelain. Fear for chipping of conventional mono glass component zirconia porcelains on the longer term and price pressure on manual application of porcelain, are possible drivers for the monolithic zirconia restorations. However, by the application of multi-glass component porcelain chipping is no longer an issue, especially with prosthetic mimetic restorations where the crown follows a model of the natural tooth in two layers: a histo-anatomic dentin layer mimicking the dentin shape of the dentition of the patient and an enamel layer. These restorations that mimic the structure of natural teeth by cognitive design of the dentin core present a new production paradigm to fabricate natural restorations of veneered zirconia using a high strength porcelain with CAD/CAM. These crowns are produced with a core of tooth-colored tetragonal zirconia, on which a high strength translucent porcelain layer has been applied and subsequently milled to size. In the subtle cooperation between the dentin-colored zirconia and the veneering porcelain, the zirconia shines through the translucent porcelain layer, all the more as the porcelain layer is thinner. This creates the natural color dynamics with color "from the inside" as found in natural elements, instead of color "on the outside", with monolithic zirconia. As a result, the natural tooth, in terms of esthetics and hardness, is approached closer than crowns made from solid monolithic zirconia. This implies that the histo-anatomic dentin core is the key to esthetic crowns.
Zirconia is the hardest known ceramic in industry and the strongest material used in dentistry, it has to be fabricated using a CAD/CAM process but not the conventional manual dental technology. Because of this monolithic zirconia does not wear itself as the normal vertical wear of 25-75 microns of natural enamel and porcelain, there are no clinical data on the fact whether as a consequence too high zirconia crowns will damage opposing dentition on the longer term. Although in two body wear testing of monolithic, veneered and glazed zirconia and their corresponding enamel antagonists showed similar wear, at least twice as much extensive, and branched enamel microcracks were observed in the samples opposing monolithic zirconia.
Monolithic zirconia
Monolithic zirconia crowns tend to be opaque in appearance with a high value and they lack translucency and fluorescence. For the sake of appearance, many dentists will not use monolithic crowns on anterior (front) teeth. Monolithic zirconia crowns are produced from a color- and structure-graded zirconia block and coated with a thin layer of glaze stains which also provides some kind of fluorescence. The "graded" zirconia crown has a darker cervical area consisting of tetragonal zirconia, a main tooth color in the buccal area, and a translucent incisal edge consisting of cubic zirconia. The only thing a dental technician has to do is to use the proper height of the zirconia block so that the crown fits in all the different color zones. Although on the outside the color gradient mimics natural teeth, they are still far from the optical, physical, biomimetic and esthetic properties of natural teeth.
To a large extent, materials selection in dentistry determine the strength and appearance of a crown. Some monolithic zirconia materials produce the strongest crowns in dentistry (the registered strength for some zirconia crown materials is near 1200 MPa), but these crowns are not usually considered to be natural enough for use in the front of the mouth. Although not as strong, some of the newer zirconia materials are better in appearance but generally still not as good as porcelain-fused crowns. By contrast, when porcelain is fused to glass-infiltrated alumina, crowns are very natural-looking and very strong, though not as strong as monolithic zirconia crowns.
Zirconia crowns are said to be less abrasive to opposing teeth than metal-ceramic crowns.
Other crown material properties to be considered are thermal conductivity and radiolucency. Stability/looseness of fit on the prepared tooth and cement gap at the margin are sometimes related to materials selection, though these crown properties are also commonly related to system and fabricating procedures.
Lithium-disilicate
Another monolithic material, lithium disilicate, produces extremely translucent leucite-reinforced crowns that often appear to be too gray in the mouth. To overcome this, the light shade polyvalent colorants take on a distinctly unnatural, bright white appearance.
However, research has shown that the shade of the supporting tooth or abutment significantly influences the final appearance of lithium disilicate crowns. Discolored foundations can adversely affect the esthetics of the restoration. Utilizing opaque ceramic cores or liners has been suggested to mitigate this issue, improving the color match and overall appearance.
Metal-ceramic crowns (P-F-M Crown)
These are a hybrid of metal and ceramic crowns. The metal part is normally made of a base metal alloy (termed bonding alloy). The properties of the metal alloy chosen should match and complement that of the ceramic to be bonded otherwise problems like delamination or fracturing of the ceramic can occur. To obtain an aesthetic finish which is able to be functional with normal mastication activity, a minimal thickness of ceramic and metallic material is required, which should be planned for during tooth preparation stage.
Ceramic bonds to the metal framework by three methods:
Compression fit (via ceramic shrinkage on firing)
Micro-mechanical retention (via surface irregularities)
Chemical union (via oxide formation)
Tissue control and gingival retraction
Gingival retraction refers to the displacement of the free gingivae. For crowns with margins which are supragingival, there is no need for gingival retraction, provided there is good moisture control.
For crown preparations which have subgingival margins, tissue control is necessary at the preparation stage and impression stage to ensure visibility, good moisture control and ensure enough bulk of impression material can be placed to accurately record the marginal areas.
Options available are gingival retraction cord, Magic Foam cord, and ExpaSyl.
Another method to expose the margins of a subgingival preparation is to use electrosurgery or crown lengthening surgery.
Tooth preparation
The design of a preparation for a tooth to accept a crown follows five basic principles:
Retention and resistance
Preservation of tooth structure
Structural durability
Marginal integrity
Preservation of the periodontium
Aesthetics can also play a role in planning the design.
Retention and resistance
As there are currently no biologically compatible cements which are able to hold the crown in place solely through their adhesive properties, the geometric form of the preparation is vital in providing retention and resistance to hold the crown in place. Within the context of prosthodontics, retention refers to resistance of movement of a restoration along the path of insertion or along the long axis of the tooth. Resistance refers to the resistance of movement of the crown by forces applied apically or in an oblique direction which prevents movement under occlusal forces. Retention is determined by the relationship between opposing surfaces of the preparation (e.g. the relationship of the buccal and lingual walls).
Taper
Theoretically, the more parallel the opposing walls of a preparation, the more retention is achieved. However this is almost impossible to achieve clinically. It is standard for preparations for full coverage crowns to slightly taper or converge in an occlusal direction. This allows the preparation to be visually inspected, prevent undercuts, compensate for crown fabrication inaccuracies and allow, at the cementation stage, for excess cement to escape with the ultimate aim of optimising the seating of the crown on the preparation. Generally axial walls prepared using a long tapered high speed burs confer a 2 - 3° taper on each wall and an overall 4 - 6° taper to the preparation. As taper increases, retention decreases so taper should be kept to a minimum whilst ensuring elimination of undercuts. An overall taper of 16° is said to be clinically achievable and being able to fulfil the aforesaid requirements. Ideally, the taper should not exceed 20 degrees as will negatively impact retention.
Length
Occluso-gingival length or height of the crown preparation affects both resistance and retention. Generally, the taller the preparation, the greater the surface area is. For the crown to be retentive enough, the length of the preparation must be greater than the height formed by the arc of the cast pivoting around a point on the margin on the opposite side of the restoration. The arc is affected by the diameter of the tooth prepared, therefore the smaller the diameter, the shorter the length of the crown needs to be to resist removal. Retention of short-walled teeth with a wide diameter can be improved by placing grooves in the axial walls, which has the effect of reducing the size of the arc.
Freedom of displacement
Retention can be improved by geometrically limiting the number of paths along which the crown can be removed from the tooth presentation, with maximum retention being reached when only one path of displacement is present. Resistance can be improved by inserting components like grooves.
Preservation of tooth structure
Preparing a tooth to accept a full coverage crown is relatively destructive. The procedure can damage the pulp irreversibly, through mechanical, thermal and chemical trauma and making the pulp more susceptible to bacterial invasion. Therefore, preparations must be as conservative as possible, whilst producing a strong retentive restoration. Although it may be seen as contradictory to the previous statement, at times, sound tooth structure may need to be sacrificed in order to prevent further more substantial and uncontrolled loss of tooth structure.
Structural durability
In order to last, the crown must be made of enough material to withstand normal masticatory function and should be contained within the space created by the tooth preparation, otherwise problems may arise with aesthetics and occlusal stability (i.e. high restorations) and cause periodontal inflammation. Depending on the material used to create the crown, minimal occlusal and axial reductions are required to house the crown.
Occlusal reduction
For gold alloys there should be 1.5mm clearance, whilst metal-ceramic crowns and full ceramic crowns require 2.0 mm. The occlusal clearance should follow the natural outline of the tooth; otherwise there may be areas of the restorations where the material may be too thin.
Functional cusp bevel
For posterior teeth, a wide bevel is required on the functional cusps, palatal cusps for maxillary teeth and buccal cusps for mandibular teeth. If this functional cusp bevel is not present and the crown is cast to replicate the correct size of the tooth, bulk of material may be too little at this point to withstand occlusal surfaces.
Axial reduction
This should allow enough thickness for the material chosen. Depending on the type of crown to be fitted, there is a minimum preparation thickness. Generally, full metal crowns require at least 0.5mm, whist metal-ceramic and full ceramic crowns require at least 1.2mm
Marginal integrity
In order for the cast restoration to last in the oral environment and to protect the underlying tooth structure, the margins between cast and tooth preparation need to be as closely adapted. The marginal line design and position should facilitate plaque control, allow for adequate thickness of the restorative material chosen therefore providing enough strength for the crown at the margin. Several types of finish line configurations have been advocated, each having some advantages and disadvantages (see the table below). Chamfer finish are normally advocated for full metal margins and shoulders are generally required to provide enough bulk for metal-ceramic crowns and full ceramic crown margins. Some evidence suggests adding a bevel to margins, especially where these are heavy, to decrease the distance between the crown and the tooth tissue.
Preservation of the periodontium
Linked to marginal integrity, placement of the finish line can directly affect the ease of manufacturing the crown and health of the periodontium. Best results are achieved where the finish line is above the gum line as this is fully cleanable. They should also be placed on enamel as this creates a better seal. Where circumstances require the margins to be below the gum line, caution is required as several problems can arise. First, there might be issues in terms of capturing the margin when making impressions during the manufacturing process leading to inaccuracies. Secondly, the biologic width, the mandatory distance (roughly 2 mm) to be left between the height of the alveolar bone and the margin of the restoration; if this distance is violated, it can result in gingival inflammation with pocket formation, gingival recession and loss of alveolar bone crest height. In these cases, crown lengthening surgery should be considered.
Special considerations
Ferrule effect
Endodontically treated teeth, especially those with little sound tooth tissue, are prone to fractures. The successful clinical outcome for these teeth relies not only on adequate root canal treatment, but also on the type of restorative treatment used, including the use of a post and core system and the type of extra-coronal restoration selected. Some evidence advocates the use of a ferrule to optimise the bio-mechanical behaviour of root-filled teeth, especially where a post and core system needs to be used.
In dentistry, the ferrule effect is, as defined by Sorensen and Engelman (1990), a "360° metal collar of the crown surrounding the parallel walls of the dentine extending coronal to the shoulder of the preparation". Like the ferrule of a pencil which encircles the junction between the rubber and the pencil shaft, the ferrule effect is believed to minimise the concentration of stresses at the junction of post and core, ultimately providing a protective effect against fractures. It also reduces stress transmission to the root due to non-axial forces applied by the post during placement or during normal function. The ferrule can also help preserve the hermetic seal of the luting cement. It has been suggested that protection acquired by the use of a ferrule occurs due to the ferrule resisting functional lever forces, wedging effect of tapered posts and lateral forces during post insertion. To make full use of the ferrule effect, the preparation needs to allow for a continuous band of dentine which should be at least 2 mm in height from the level of the preparation margin and with the band being at least 1 mm in thickness.
It has been shown, however, that whilst the absence of a 360° ferrule can increase the risk of fracture of root-filled teeth restored with fiber post and cores and crowns, having insufficient coronal walls poses an even greater one.
Stainless steel crowns for posterior primary dentition
Stainless steel preformed metal crowns are the treatment of choice for the restoration of posterior primary teeth. A systemic review found that it has the highest success rate (96.1%). In order to accept a stainless steel crown, the entire occlusal surface should be reduced by 1–1.5 mm and interproximally contacts should be cleared by cutting a thin mesial and distal portion or slice subgingivally by holding the tip of a thin high-speed bur at 15–20° relative to the long axis of the tooth, to avoid the creation of a shoulder. No preparation of the buccal or lingual/palatal surfaces is required. Stainless steel crowns can be made esthetic by veneering composite using the open face technique or composite veneering done after sand-blasting SSCs. Also, composite veneering can be done after preparing retentive grooves on the buccal surface of stainless steel crowns.
Hall technique
The Hall technique is a non-invasive treatment for decayed posterior primary teeth where caries are sealed under a preformed stainless steel crown. This technique requires no tooth preparation.
Construction and fit of temporary crown restorations
It is very likely that once a tooth has been prepared and whilst waiting for the definitive restoration, the tooth preparation is fitted with a temporary crown.
Need for temporary restorations
Temporisation is important after tooth preparation in order to:
Protect from and prevent bacterial invasion of newly exposed dentinal tubules, leading to pulpal inflammation and necrosis;
Prevent gingival growth in the area created by the tooth preparation;
Allow area to be cleaned more effectively, decreasing the incidence of bleeding and gingival inflammation at the time of fitting definitive restoration;
Maintain occlusal and approximal contacts therefore preventing over-eruption, rotation and closing of spaces;
Aesthetic reasons;
Temporary crowns can also play a diagnostic role in treatment planning where there is a need for occlusal, aesthetic or periodontal changes.
Types of temporary crowns
Temporary crowns can be described by:
The expected or planned duration of temporisation:
Short term
Medium term
Long term
The way or the place the temporary restoration is made:
Direct or chair-side
Indirect or laboratory-made
The aesthetics or look of the material of construction
Metal
Cast
Preformed
Tooth coloured
Plastic pre-formed (e.g. polycarboxylate and acrylic)
Resin composites
Duration of temporisation
Temporary crowns can be described as short-term, if used for a few days, medium-term, if their planned use for several weeks and long-term if their planned use is for several months. The choice in length of temporisation often relates to the complexity of restorative work planned. Short-term temporary crowns are generally appropriate for simple restorative cases whilst complex cases involving more than one tooth often require long-term temporary crowns.
Direct vs. indirect restorations
Temporary crowns can either be direct, if constructed by the dentist in the clinic, or indirect if they are made off-site, usually in a dental laboratory. Generally direct temporary crowns tend to be for short-term use. Where medium-term or long-term temporisation is required, the use of indirect temporary crowns should be considered.
Temporary crown materials
There are several materials that can be used to construct temporary crowns. Direct temporary crowns are either made using metal or plastic pre-formed crowns, chemically cured or light-cured resins or resin composites. Indirect restorations are either made of chemically cured acrylic, heat-cured acrylic or cast in metal.
Cementation of temporary crowns
The purpose of temporary luting agents is to fill the space between the crown preparation and the temporary restoration. Unlike cementation of definitive crowns, temporary crowns should be relatively easy to remove. Adhesive cements should not be used and softer cements are preferred to allow for the easy removal of both temporary cements and crowns. This is crucial as remnants of temporary cement left on the tooth surface can compromise gingival health and interfere with accurate seating of the final restoration and permanent cement attachment. Provisional cements should also be strong enough to avoid being deformed or fractured during the provisional period.
Zinc Oxide Eugenol (ZOE) temporary luting cements
These are commonly used because of their low tensile strength and lack of adhesion which provides ease of removal. These products should not be used when resin composite is to be planned for bonding the definitive crown as eugenol is able to infiltrate and diffuse through dentine; contaminating tooth surface and compromising bonding by inhibiting polymerisation of resin. Commercially available products include RelyX Temp E (3M ESPE), Temp-Bond (Kerr) and Flow Temp (Premier Dental Products).
Non-eugenol temporary luting cements
Non-eugenol cements replace eugenol with several types of carboxylic acids which do not inhibit definitive cementation. These cements are compatible with temporary resin materials and definitive resin cements and have increased retention when compared to ZOE containing cements. Examples of commercially available products include RelyX Temp NE (3M ESPE) and Temp-Bond NE (Kerr).
Polycarboxylate temporary luting cements
This hydrophilic cement has the benefit of minimal effects on temporary resin containing agents and weak adhesion to tooth tissue which increases ease of removal. This cement is the easiest to clean out of all the provisional cement types. Examples include Ultradent and Hy-Bond (Shofu Dental).
Resin temporary luting cements
The advantages of these cements include superior aesthetics, greater strength, superb retention and ease of cleaning. However, amongst the drawbacks of this cement is the higher rate of discolouration, microleakage and odour experienced. Commercially available examples of temporary resins cements include Systemp.link (Ivoclar Vivadent), Temp-Bond Clear (Kerr) and ImProv (Nobel Biocare).
Tooth preparation impressions
Once the tooth in question has been prepared with acceptable dimensions, it is equally important to make an accurate and dimensionally stable record or impression of the preparation or dental implant, surrounding hard and soft tissues as well as the opposing dental arch so that the restoration created will conform to the required dimensions and ensure the fit is as close as possible without having to make many modifications chair-side.
Impressions can be made digitally or by conventional technique. With regards to conventional impression techniques, the materials selected should have appropriate physical properties and handling characteristics to allow enough detail reproduction and durability when casting a model, including the ability to withstand effective decontamination procedures. Generally, impressions of the arch where the preparation is made are in addition silicone using the "wash impression" technique; impressions of the opposing arch are made in alginate.
Digital impressions can be made using dedicated optical scanners. A review suggests that digital impressions provide the same accuracy as conventional impressions and are found to be more comfortable for patients and easier for dental practitioners.
Crown manufacture using CAD/CAM
CEREC
Chairside CAD/CAM dentistry
The CAD/CAM method of fabricating all-ceramic restorations is by electronically capturing and storing a photographic image of the prepared tooth and, using computer technology, crafting a 3D restoration design that conforms to all the necessary specifications of the proposed inlay, onlay or single-unit crown; there is no impression. After selecting the proper features and making various decisions on the computerized model, the dentist directs the computer to send the information to a local milling machine. This machine will then use its specially designed diamond burs to mill the restoration from a solid ingot of a ceramic of pre-determined shade to match the patient's tooth. After about 20 minutes, the restoration is complete, and the dentist sections it from the remainder of the unmilled ingot and tries it in the mouth. If the restoration fits well, the dentist can cement the restoration immediately. A dental CAD/CAM machine costs roughly $100,000, with continued purchase of ceramic ingots and milling burs. Because of high costs, the usual and customary fee for making a CAD/CAM crown in the dentist's office is often slightly higher than having the same crown made in a dental laboratory.
Typically, over 95% of the restorations made using dental CAD/CAM and Vita Mark I and Mark II blocks are still clinically successful after five years. Further, at least 90% of restorations still function successfully after 10 years. Advantages of the Mark II blocks over ceramic blocks include: they wear down as fast as natural teeth, their failure loads are very similar to those of natural teeth, and the wear pattern of Mark II against enamel is similar to that of enamel against enamel.
In recent years, the technological advances afforded by CAD/CAM dentistry offer viable alternatives to the traditional crown restoration in many cases. Where the traditional indirectly fabricated crown requires a tremendous amount of surface area to retain the normal crown, potentially resulting in the loss of healthy, natural tooth structure for this purpose, the all-porcelain CAD/CAM crown can be predictably used with significantly less surface area. As a matter of fact, the more enamel that is retained, the greater the likelihood of a successful outcome. As long as the thickness of porcelain on the top, chewing portion of the crown is 1.5mm thick or greater, the restoration can be expected to be successful. The side walls which are normally totally sacrificed in the traditional crown are generally left far more intact with the CAD/CAM option. In regards to post and core buildups, these are generally contraindicated in CAD/CAM crowns as the resin bonding materials do best bonding the etched porcelain interface to the etched enamel/dentin interfaces of the natural tooth itself. The crownlay is also an excellent alternative to the post and core buildup when restoring a root canal-treated tooth.
Crown removal
At times it may be necessary to remove crown restorations to enable treatment of the tooth tissue underneath, especially to enable for non-surgical endodontic treatment of a necrotic or previously treated pulp. Several methods are available and the choice is guided normally by the nature and quality of the crown restoration., i.e., whether it is to be retained or to be replaced.
Factors to consider when deciding whether to retain or remove the crown include:
Replacement cost (time and financial)
Aesthetic
Ease of removal
Marginal integrity
Planned restoration (including change from a crown to a bridge, or adapting the crown design to act as abutment for a partial denture)
Access required to treat the tooth safely and effectively (especially with regards to access cavity design)
Temporary crowns are easy to remove and replace therefore do not pose problems.
Before removing definitive crown restorations it is important to plan for a provisional crown, especially if the crown to be removed is expected to be damaged in the process. This usually involves making an impression of the crown so a temporary can be fabricated chair-side or made by the dental laboratory. Several tools and methods are available, which can be classified by how conservative they are to the crown. Normally the tooth, if heavily damaged, should be restored prior to a new crown (whether temporary or definitive) is fitted.
Matrix bands
Application of a matrix band which is burnished into the undercuts and pulled vertically.
Ultrasonic
An ultrasonic tip can be applied to a cast metal crown to disrupt the cement lute. This method should be avoided with ceramic restorations as this may lead to fractures.
Forceps and crown tractors
Crown tractors and forceps can be used to grip the restoration and dislodge it from the tooth preparation. Crown tractors are designed to have rubber grips and powder on their beaks to reduce the risk of damaging ceramic restorations. Crown tractors are quite effective in removing crowns luted with temporary cements.
Sticky sweet method or Richwill crown and bridge remover
A thermoplastic pliable resin is softened in warm water then placed on the occlusal surface of the crown to be removed. The patient is then asked to bite down, compressing the resin block to two-thirds of its original thickness. The patient is then asked to open the mouth quickly, which should generate enough force to displace the restoration. This method however, is not very effective and has a risk of damaging restorations on or accidentally extracting the opposing tooth. Therefore, before using this method, it is important to look at the state of the opposing tooth.
Tapping and pneumatic tools
Sliding hammers work by using a tip to engage the crown margin and sliding the weight along the shaft and tapping this to loosen the restoration. Several versions are available. Some are weighted, others are spring loaded. This system is uncomfortable for the patient and is not always successful. It is also contraindicated for periodontically involved teeth, as it can cause unwanted extractions. This system can also damage the ceramic margins.
Wedging devices
A slot is cut using the side of a long tapered bur, normally on the buccal surface of the crown, all the way into the cement lute. A flat plastic instrument, straight Warrick James, Couplands elevators or dedicated systems such as the WamKey, is inserted into the slot created to wedge the crown apart from the tooth.
Metalift crown and bridge removal system
Based on the "jack-screw" principle, the Metalift system works by drilling a precision channel through the occlusal surface of a cast restoration, then with a special bur, the area around the periphery of the hole is undermined before a threaded screw is wound into the space. As the screw comes in contact with the core of the restoration, the continued rotation of the screw results in a jacking force that displaces the crown from the preparation. This system can be used to remove both all metal crowns and metal-ceramic crowns, although with metal-ceramic crowns care should be taken to remove enough ceramic from the area where the hole created to reduce the chances of fracture. The minimum thickness of metal required for the lifting action is approximately 0.5mm. The damage can repaired with a plastic filling material.
Burs
The crown can be simply sectioned using a bur.
History
There is evidence of gold dental prosthesis dating back to the Etruscans. Ancient Roman dentists also adopted these tools.
| Biology and health sciences | Dental treatments | Health |
1343446 | https://en.wikipedia.org/wiki/Pandanus | Pandanus | Pandanus is a genus of monocots with about 578 accepted species. They are palm-like, dioecious trees and shrubs native to the Old World tropics and subtropics. Common names include pandan, screw palm and screw pine. They are classified in the order Pandanales, family Pandanaceae.
Description
The species vary in size from small shrubs less than tall, to medium-sized trees tall, typically with a broad canopy, heavy fruit, and moderate growth rate. The trunk is stout, wide-branching, and ringed with many leaf scars. Mature plants can have branches. Depending on the species, the trunk can be smooth, rough, or warty. The roots form a pyramidal tract to hold the trunk. They commonly have many thick stilt roots near the base, which provide support as the tree grows top-heavy with leaves, fruit, and branches. These roots are adventitious and often branched. The top of the plant has one or more crowns of strap-shaped leaves that may be spiny, varying between species from to or longer, and from up to broad.
They are dioecious, with male and female flowers produced on different plants. The flowers of the male tree are long and fragrant, surrounded by narrow, white bracts. The female tree produces flowers with round fruits that are also bract-surrounded. The individual fruit is a drupe, and these merge to varying degrees forming multiple fruit, a globule structure, in diameter and have many prism-like sections, resembling the fruit of the pineapple. Typically, the fruit changes from green to bright orange or red as it matures. The fruits can stay on the tree for more than 12 months.
Taxonomy
Often called pandanus palms, these plants are not closely related to palm trees. The genus is named after the Malay word pandan given to Pandanus amaryllifolius, the genus's most commonly known species. The name is derived from Proto-Austronesian *paŋudaN (which became Proto-Malayo-Polynesian *pangdan and Proto-Oceanic *padran). It has many cognates in Austronesian languages, underscoring its importance in Austronesian cultures, including Atayal pangran; Kavalan pangzan; Thao panadan; Tagalog pandan; Chamorro pahong; Manggarai pandang; Malagasy fandrana, Tongan fā; Tahitian fara; Hawaiian hala all referring to plants of similar characteristics and/or uses whether in the same genus (particularly Pandanus tectorius) or otherwise (in the case of Māori whara or hara; e.g. harakeke).
The oldest fossil of the genus is Pandanus estellae which is known from a silicified fruit found in Queensland, Australia, dating to the Oligocene epoch around 32–28 million years ago.
Selected species
Note: several species previously placed in Pandanus subgenus Acrostigma are now in the distinct genus Benstonea.
Pandanus aldabraensis H.St.John
Pandanus amaryllifolius Roxb. ex Lindl.
Pandanus balfourii Martelli
Pandanus barkleyi Balf.f.
Pandanus boninensis Warb.
Pandanus candelabrum P.Beauv.
Pandanus carmichaelii R.E.Vaughan & Wiehe
Pandanus ceylanicus Solms
Pandanus christmatensis Martelli
Pandanus clandestinus Stone
Pandanus conglomeratus Balf.f.
Pandanus conoideus Lam.
Pandanus decastigma B.C.Stone
Pandanus decipiens Martelli
Pandanus decumbens Solms
Pandanus drupaceus Thouars
Pandanus elatus Ridl.
Pandanus eydouxia Balf.f.
Pandanus fanningensis H.St.John
Pandanus forsteri C.Moore & F.Muell.
Pandanus furcatus Roxb.
Pandanus gabonensis Huynh
Pandanus glaucocephalus R.E.Vaughan & Wiehe
Pandanus grayorum
Pandanus halleorum B.C.Stone
Pandanus heterocarpus Balf.f.
Pandanus iceryi Horne ex Balf.f.
Pandanus incertus R.E.Vaughan & Wiehe
Pandanus joskei Horne ex Balf.f.
Pandanus julianettii Martelli
Pandanus kaida Kurz
Pandanus kajui Beentje
Pandanus lacuum H.St.John ex B.C.Stone
Pandanus laxespicatus Martelli
Pandanus livingstonianus Rendle
Pandanus leram
Pandanus microcarpus Balf.f.
Pandanus montanus Bory
Pandanus multispicatus Balf.f.
Pandanus odorifer (Forssk.) Kuntze
Pandanus obeliscus
Pandanus palustris Thouars
Pandanus parvicentralis Huynh
Pandanus prostratus Balf.f.
Pandanus pyramidalis Barkly ex Balf.f.
Pandanus rigidifolius R.E.Vaughan & Wiehe
Pandanus sechellarum Balf.f.
Pandanus spathulatus Martelli
Pandanus spiralis R.Br.
Pandanus tectorius Parkinson ex Du Roi
Pandanus tenuifolius Balf.f.
Pandanus teuszii Warb.
Pandanus thomensis Henriq.
Pandanus tonkinensis B.C.Stone
Pandanus utilis Bory
Pandanus vandermeeschii Balf.f.
Pandanus verecundus Stone
Distribution and habitat
The greatest number of species are found in Madagascar and Malaysia.
Ecology
These plants grow from sea level to an altitude of . Pandanus trees are of cultural, health, and economic importance in the Pacific, second only to the coconut on atolls. They grow wild mainly in semi-natural vegetation in littoral habitats throughout the tropical and subtropical Pacific, where they can withstand drought, strong winds, and salt spray. They propagate readily from seed, but popular cultivars are also widely propagated from branch cuttings by local people.
Species growing on exposed coastal headlands and along beaches have thick 'stilt roots' as anchors in the loose sand. Those stilt roots emerge from the stem, usually close to but above the ground, which helps to keep the plants upright and secure them to the ground.
While pandanus are distributed throughout the tropical and subtropical islands and coastlines of the Atlantic, Indian and Pacific Oceans, they are most numerous on the low islands and barren atolls of Polynesia and Micronesia. Other species are adapted to mountain habitats and riverine forests.
The tree is grown and propagated from shoots that form spontaneously in the axils of lower leaves. Pandanus fruits are eaten by animals including bats, rats, crabs, and elephants, but the vast majority of species are dispersed primarily by water. Its fruit can float and spread to other islands without help from humans.
Uses
Pandanus has multiple uses, which is dependent in part on each type and location. Some pandanus are a source of food, while others provide raw material for clothing, basket weaving and shelter.
Pandanus leaves are used for handicrafts. Artisans collect the leaves from plants in the wild, cutting only mature leaves so that the plant will naturally regenerate. The leaves are sliced into fine strips and sorted for further processing. Weavers produce basic pandan mats of standard size or roll the leaves into pandan ropes for other designs. This is followed by a coloring process, in which pandan mats are placed in drums with water-based colors. After drying, the colored mats are shaped into final products, such as placemats or jewelry boxes. Final color touch-ups may be applied. The species in Hawaiʻi are called hala, and only the dry leaves (lauhala) are collected and used for Lauhala weaving.
Pandanus leaves from Pandanus amaryllifolius are used widely in Southeast Asian and South Asian cuisines to add a distinct aroma to various dishes and to complement flavors like chocolate. Because of their similarity in usage, pandan leaves are sometimes referred to as the "vanilla of Asia." Fresh leaves are typically torn into strips, tied in a knot to facilitate removal, placed in the cooking liquid, then removed at the end of cooking. Dried leaves and bottled extract may be bought in some places. Finely sliced pandan leaves are used as fragrant confetti for Malay weddings, graves etc.
Pandan leaves are known as Daun pandan in Indonesian and Malaysian Malay; Dahon ng pandan (lit. "pandan leaf") or simply pandan in Filipino; 斑蘭 (bān lán) in Mandarin; as ใบเตย (bai toei; ) in Thai, lá dứa in Vietnamese; pulao data in Bengali; and rampe in Sinhalese and Hindi.
In India, particularly in Nicobar Islands, pandanus fruit is staple food of Shompen people and Nicobarese people.
In Sri Lanka, pandan leaves are used heavily in both vegetable and meat dishes and are often grown in homes. It is common practice to add a few pieces of pandan leaf when cooking red or white rice as well.
In Southeast Asia, pandan leaves are mainly used in sweets such as coconut jam and pandan cake. In Indonesia and Malaysia, pandan is also added to rice and curry dishes such as nasi lemak. In the Philippines, pandan leaves are commonly paired with coconut meat (a combination referred to as buko pandan) in various desserts and drinks like maja blanca and gulaman.
In Indian cooking, the leaf is added whole to biryani, a kind of rice pilaf, made with ordinary rice (as opposed to that made with the premium-grade basmati rice). The basis for this use is that both basmati and pandan leaf contains the same aromatic flavoring ingredient, 2-acetyl-1-pyrroline. In Sri Lanka, pandan leaves are a major ingredient used in the country's cuisine.
Kewra (also spelled Kevda or Kevada) is an extract distilled from the pandan flower, used to flavor drinks and desserts in Indian cuisine. Also, kewra or kevada is used in religious worship, and the leaves are used to make hair ornaments worn for their fragrance as well as decorative purpose in western India.
Species with large and medium fruit are edible, notably the many cultivated forms of P. tectorius (P. pulposus) and P. utilis. The ripe fruit can be eaten raw or cooked, while partly ripe fruit should be cooked first. Small-fruited pandanus may be bitter and astringent.
Karuka nuts (P. julianettii) are an important staple food in New Guinea. Over 45 cultivated varieties are known. Entire households will move, and in some areas will speak a pandanus language at harvest time. The taste is like coconut or walnuts.
Throughout Oceania, almost every part of the plant is used, with various species different from those used in Southeast Asian cooking. Pandanus trees provide materials for housing; clothing and textiles including the manufacture of dilly bags (carrying bags), fine mats or ie toga; sails, food, medication, decorations, fishing, and religious uses. In the Vanuatu Archipelago, natives make woven fish traps from the hard interior root of the pandanus, made like a cage having a narrow entrance.
| Biology and health sciences | Monocots | null |
1343550 | https://en.wikipedia.org/wiki/P%20wave | P wave | A P wave (primary wave or pressure wave) is one of the two main types of elastic body waves, called seismic waves in seismology. P waves travel faster than other seismic waves and hence are the first signal from an earthquake to arrive at any affected location or at a seismograph. P waves may be transmitted through gases, liquids, or solids.
Nomenclature
The name P wave can stand for either pressure wave (as it is formed from alternating compressions and rarefactions) or primary wave (as it has high velocity and is therefore the first wave to be recorded by a seismograph). The name S wave represents another seismic wave propagation mode, standing for secondary or shear wave, a usually more destructive wave than the primary wave.
Seismic waves in the Earth
Primary and secondary waves are body waves that travel within the Earth. The motion and behavior of both P and S waves in the Earth are monitored to probe the interior structure of the Earth. Discontinuities in velocity as a function of depth are indicative of changes in phase or composition. Differences in arrival times of waves originating in a seismic event like an earthquake as a result of waves taking different paths allow mapping of the Earth's inner structure.
P wave shadow zone
Almost all the information available on the structure of the Earth's deep interior is derived from observations of the travel times, reflections, refractions and phase transitions of seismic body waves, or normal modes. P waves travel through the fluid layers of the Earth's interior, and yet they are refracted slightly when they pass through the transition between the semisolid mantle and the liquid outer core. As a result, there is a P wave "shadow zone" between 103° and 142° from the earthquake's focus, where the initial P waves are not registered on seismometers. In contrast, S waves do not travel through liquids.
As an earthquake warning
Advance earthquake warning is possible by detecting the nondestructive primary waves that travel more quickly through the Earth's crust than do the destructive secondary and Rayleigh waves.
The amount of warning depends on the delay between the arrival of the P wave and other destructive waves, generally on the order of seconds up to about 60 to 90 seconds for deep, distant, large quakes such as the 2011 Tohoku earthquake. The effectiveness of a warning depends on accurate detection of the P waves and rejection of ground vibrations caused by local activity (such as trucks or construction). Earthquake early warning systems can be automated to allow for immediate safety actions, such as issuing alerts, stopping elevators at the nearest floors, and switching off utilities.
Propagation
Velocity
In isotropic and homogeneous solids, a P wave travels in a straight line longitudinally; thus, the particles in the solid vibrate along the axis of propagation (the direction of motion) of the wave energy. The velocity of P waves in that kind of medium is given by
where is the bulk modulus (the modulus of incompressibility), is the shear modulus (modulus of rigidity, sometimes denoted as and also called the second Lamé parameter), is the density of the material through which the wave propagates, and is the first Lamé parameter.
In typical situations in the interior of the Earth, the density usually varies much less than or , so the velocity is mostly "controlled" by these two parameters.
The elastic moduli P wave modulus, , is defined so that and thereby
Typical values for P wave velocity in earthquakes are in the range 5 to 8 km/s. The precise speed varies according to the region of the Earth's interior, from less than 6 km/s in the Earth's crust to 13.5 km/s in the lower mantle, and 11 km/s through the inner core.
Geologist Francis Birch discovered a relationship between the velocity of P waves and the density of the material the waves are traveling in:
which later became known as Birch's law. (The symbol is an empirically tabulated function, and is a constant.)
| Physical sciences | Seismology | Earth science |
30125638 | https://en.wikipedia.org/wiki/Cubic%20metre | Cubic metre | The cubic metre (in Commonwealth English and international spelling as used by the International Bureau of Weights and Measures) or cubic meter (in American English) is the unit of volume in the International System of Units (SI). Its symbol is m3. It is the volume of a cube with edges one metre in length. An alternative name, which allowed a different usage with metric prefixes, was the stère, still sometimes used for dry measure (for instance, in reference to wood). Another alternative name, no longer widely used, was the kilolitre.
Conversions
{|
|-
|rowspan=6 valign=top|1 cubic metre
|= litres (exactly)
|-
|≈ 35.3 cubic feet
|-
|≈ 1.31 cubic yards
|-
|≈ 6.29 oil barrels
|-
|≈ 220 imperial gallons
|-
|≈ 264 US fluid gallons
|}
A cubic metre of pure water at the temperature of maximum density (3.98 °C) and standard atmospheric pressure (101.325 kPa) has a mass of , or one tonne. At 0 °C, the freezing point of water, a cubic metre of water has slightly less mass, 999.972 kilograms.
A cubic metre is sometimes abbreviated to , , , , , , when superscript characters or markup cannot be used (e.g. in some typewritten documents and postings in Usenet newsgroups). The "cubic metre" symbol is encoded by Unicode at code point .
Multiples and submultiples
Multiples
Cubic decametre
the volume of a cube of side length one decametre (10 m)
equal to a megalitre
1 dam3 = = 1 ML
Cubic hectometre
the volume of a cube of side length one hectometre (100 m)
equal to a gigalitre
in civil engineering abbreviated MCM for million cubic metres
1 hm3 = = 1 GL
Cubic kilometre
the volume of a cube of side length one kilometre ()
equal to a teralitre
1 km3 = = 1 TL (810713.19 acre-feet; 0.239913 cubic miles)
Submultiples
Cubic decimetre
the volume of a cube of side length one decimetre (0.1 m)
equal to a litre
1 dm3 = 0.001 m3 = 1 L
(also known as DCM (=Deci Cubic Meter) in Rubber compound processing)
Cubic centimetre
the volume of a cube of side length one centimetre (0.01 m)
equal to a millilitre
1 cm3 = = 10−6 m3 = 1 mL
Cubic millimetre
the volume of a cube of side length one millimetre (0.001 m)
equal to a microlitre
1 mm3 = = 10−9 m3 = 1 μL
| Physical sciences | Volume | Basics and measurement |
28691804 | https://en.wikipedia.org/wiki/Mineral%20physics | Mineral physics | Mineral physics is the science of materials that compose the interior of planets, particularly the Earth. It overlaps with petrophysics, which focuses on whole-rock properties. It provides information that allows interpretation of surface measurements of seismic waves, gravity anomalies, geomagnetic fields and electromagnetic fields in terms of properties in the deep interior of the Earth. This information can be used to provide insights into plate tectonics, mantle convection, the geodynamo and related phenomena.
Laboratory work in mineral physics require high pressure measurements. The most common tool is a diamond anvil cell, which uses diamonds to put a small sample under pressure that can approach the conditions in the Earth's interior.
Creating high pressures
Shock compression
Many of the pioneering studies in mineral physics involved explosions or projectiles that subject a sample to a shock. For a brief time interval, the sample is under pressure as the shock wave passes through. Pressures as high as any in the Earth have been achieved by this method. However, the method has some disadvantages. The pressure is very non-uniform and is not adiabatic, so the pressure wave heats the sample up in passing. The conditions of the experiment must be interpreted in terms of a set of pressure-density curves called Hugoniot curves.
Multi-anvil press
Multi-anvil presses involve an arrangement of anvils to concentrate pressure from a press onto a sample. Typically the apparatus uses an arrangement eight cube-shaped tungsten carbide anvils to compress a ceramic octahedron containing the sample and a ceramic or Re metal furnace. The anvils are typically placed in a large hydraulic press. The method was developed by Kawai and Endo in Japan. Unlike shock compression, the pressure exerted is steady, and the sample can be heated using a furnace. Pressures of about 28 GPa (equivalent to depths of 840 km), and temperatures above 2300 °C, can be attained using WC anvils and a lanthanum chromite furnace. The apparatus is very bulky and cannot achieve pressures like those in the diamond anvil cell (below), but it can handle much larger samples that can be quenched and examined after the experiment. Recently, sintered diamond anvils have been developed for this type of press that can reach pressures of 90 GPa (2700 km depth).
Diamond anvil cell
The diamond anvil cell is a small table-top device for concentrating pressure. It can compress a small (sub-millimeter sized) piece of material to extreme pressures, which can exceed 3,000,000 atmospheres (300 gigapascals). This is beyond the pressures at the center of the Earth. The concentration of pressure at the tip of the diamonds is possible because of their hardness, while their transparency and high thermal conductivity allow a variety of probes can be used to examine the state of the sample. The sample can be heated to thousands of degrees.
Creating high temperatures
Achieving temperatures found within the interior of the earth is just as important to the study of mineral physics as creating high pressures. Several methods are used to reach these temperatures and measure them. Resistive heating is the most common and simplest to measure. The application of a voltage to a wire heats the wire and surrounding area. A large variety of heater designs are available including those that heat the entire diamond anvil cell (DAC) body and those that fit inside the body to heat the sample chamber. Temperatures below 700 °C can be reached in air due to the oxidation of diamond above this temperature. With an argon atmosphere, higher temperatures up to 1700 °C can be reached without damaging the diamonds. A tungsten resistive heater with Ar in a BX90 DAC was reported to achieve temperatures of 1400 °C.
Laser heating is done in a diamond-anvil cell with Nd:YAG or lasers to achieve temperatures above 6000k. Spectroscopy is used to measure black-body radiation from the sample to determine the temperature. Laser heating is continuing to extend the temperature range that can be reached in diamond-anvil cell but suffers two significant drawbacks. First, temperatures below 1200 °C are difficult to measure using this method. Second, large temperature gradients exist in the sample because only the portion of sample hit by the laser is heated.
Properties of materials
Equations of state
To deduce the properties of minerals in the deep Earth, it is necessary to know how their density varies with pressure and temperature. Such a relation is called an equation of state (EOS). A simple example of an EOS that is predicted by the Debye model for harmonic lattice vibrations is the Mie-Grünheisen equation of state:
where is the heat capacity and is the Debye gamma. The latter is one of many Grünheisen parameters that play an important role in high-pressure physics. A more realistic EOS is the Birch–Murnaghan equation of state.
Interpreting seismic velocities
Inversion of seismic data give profiles of seismic velocity as a function of depth. These must still be interpreted in terms of the properties of the minerals. A very useful heuristic was discovered by Francis Birch: plotting data for a large number of rocks, he found a linear relation of the compressional wave velocity of rocks and minerals of a constant average atomic weight with density :
.
This relationship became known as Birch's law.
This makes it possible to extrapolate known velocities for minerals at the surface to predict velocities deeper in the Earth.
Other physical properties
Viscosity
Creep (deformation)
Melting
Electrical conduction and other transport properties
Methods of crystal interrogation
There are a number of experimental procedures designed to extract information from both single and powdered crystals. Some techniques can be used in a diamond anvil cell (DAC) or a multi anvil press (MAP). Some techniques are summarized in the following table.
First principles calculations
Using quantum mechanical numerical techniques, it is possible to achieve very accurate predictions of crystal's properties including structure, thermodynamic stability, elastic properties and transport properties. The limit of such calculations tends to be computing power, as computation run times of weeks or even months are not uncommon.
History
The field of mineral physics was not named until the 1960s, but its origins date back at least to the early 20th century and the recognition that the outer core is fluid because seismic work by Oldham and Gutenberg showed that it did not allow shear waves to propagate.
A landmark in the history of mineral physics was the publication of Density of the Earth by Erskine Williamson, a mathematical physicist, and Leason Adams, an experimentalist. Working at the Geophysical Laboratory in the Carnegie Institution of Washington, they considered a problem that had long puzzled scientists. It was known that the average density of the Earth was about twice that of the crust, but it was not known whether this was due to compression or changes in composition in the interior. Williamson and Adams assumed that deeper rock is compressed adiabatically (without releasing heat) and derived the Adams–Williamson equation, which determines the density profile from measured densities and elastic properties of rocks. They measured some of these properties using a 500-ton hydraulic press that applied pressures of up to 1.2 gigapascals (GPa). They concluded that the Earth's mantle had a different composition than the crust, perhaps ferromagnesian silicates, and the core was some combination of iron and nickel. They estimated the pressure and density at the center to be 320 GPa and 10,700 kg/m3, not far off the current estimates of 360 GPa and 13,000 kg/m3.
The experimental work at the Geophysical Laboratory benefited from the pioneering work of Percy Bridgman at Harvard University, who developed methods for high-pressure research that led to a Nobel Prize in Physics. A student of his, Francis Birch, led a program to apply high-pressure methods to geophysics.
Birch extended the Adams-Williamson equation to include the effects of temperature. In 1952, he published a classic paper, Elasticity and constitution of the Earth's interior, in which he established some basic facts: the mantle is predominantly silicates; there is a phase transition between the upper and lower mantle associated with a phase transition; and the inner and outer core are both iron alloys.
| Physical sciences | Minerals | Earth science |
28691929 | https://en.wikipedia.org/wiki/Protective%20relay | Protective relay | In electrical engineering, a protective relay is a relay device designed to trip a circuit breaker when a fault is detected. The first protective relays were electromagnetic devices, relying on coils operating on moving parts to provide detection of abnormal operating conditions such as over-current, overvoltage, reverse power flow, over-frequency, and under-frequency.
Microprocessor-based solid-state digital protection relays now emulate the original devices, as well as providing types of protection and supervision impractical with electromechanical relays. Electromechanical relays provide only rudimentary indication of the location and origin of a fault. In many cases a single microprocessor relay provides functions that would take two or more electromechanical devices. By combining several functions in one case, numerical relays also save capital cost and maintenance cost over electromechanical relays. However, due to their very long life span, tens of thousands of these "silent sentinels" are still protecting transmission lines and electrical apparatus all over the world. Important transmission lines and generators have cubicles dedicated to protection, with many individual electromechanical devices, or one or two microprocessor relays.
The theory and application of these protective devices is an important part of the education of a power engineer who specializes in power system protection. The need to act quickly to protect circuits and equipment often requires protective relays to respond and trip a breaker within a few thousandths of a second. In some instances these clearance times are prescribed in legislation or operating rules. A maintenance or testing program is used to determine the performance and availability of protection systems.
Based on the end application and applicable legislation, various standards such as ANSI C37.90, IEC255-4, IEC60255-3, and IAC govern the response time of the relay to the fault conditions that may occur.
Operation principles
Electromechanical protective relays operate by either magnetic attraction, or magnetic induction. Unlike switching type electromechanical relays with fixed and usually ill-defined operating voltage thresholds and operating times, protective relays have well-established, selectable, and adjustable time and current (or other operating parameter) operating characteristics. Protection relays may use arrays of induction disks, shaded-pole, magnets, operating and restraint coils, solenoid-type operators, telephone-relay contacts, and phase-shifting networks.
Protective relays can also be classified by the type of measurement they make. A protective relay may respond to the magnitude of a quantity such as voltage or current. Induction relays can respond to the product of two quantities in two field coils, which could for example represent the power in a circuit.
"It is not practical to make a relay that develops a torque equal to the quotient of two a.c. quantities. This, however is not important; the only significant condition for a relay is its setting and the setting can be made to correspond to a ratio regardless of the component values over a wide range."
Several operating coils can be used to provide "bias" to the relay, allowing the sensitivity of response in one circuit to be controlled by another. Various combinations of "operate torque" and "restraint torque" can be produced in the relay.
By use of a permanent magnet in the magnetic circuit, a relay can be made to respond to current in one direction differently from in another. Such polarized relays are used on direct-current circuits to detect, for example, reverse current into a generator. These relays can be made bistable, maintaining a contact closed with no coil current and requiring reverse current to reset. For AC circuits, the principle is extended with a polarizing winding connected to a reference voltage source.
Lightweight contacts make for sensitive relays that operate quickly, but small contacts can't carry or break heavy currents. Often the measuring relay will trigger auxiliary telephone-type armature relays.
In a large installation of electromechanical relays, it would be difficult to determine which device originated the signal that tripped the circuit. This information is useful to operating personnel to determine the likely cause of the fault and to prevent its re-occurrence. Relays may be fitted with a "target" or "flag" unit, which is released when the relay operates, to display a distinctive colored signal when the relay has tripped.
Types according to construction
Electromechanical
Electromechanical relays can be classified into several different types as follows:
"Armature"-type relays have a pivoted lever supported on a hinge or knife-edge pivot, which carries a moving contact. These relays may work on either alternating or direct current, but for alternating current, a shading coil on the pole is used to maintain contact force throughout the alternating current cycle. Because the air gap between the fixed coil and the moving armature becomes much smaller when the relay has operated, the current required to maintain the relay closed is much smaller than the current to first operate it. The "returning ratio" or "differential" is the measure of how much the current must be reduced to reset the relay.
A variant application of the attraction principle is the plunger-type or solenoid operator. A reed relay is another example of the attraction principle.
"Moving coil" meters use a loop of wire turns in a stationary magnet, similar to a galvanometer but with a contact lever instead of a pointer. These can be made with very high sensitivity. Another type of moving coil suspends the coil from two conductive ligaments, allowing very long travel of the coil.
Induction disc overcurrent relay
"Induction" disk meters work by inducing currents in a disk that is free to rotate; the rotary motion of the disk operates a contact. Induction relays require alternating current; if two or more coils are used, they must be at the same frequency otherwise no net operating force is produced. These electromagnetic relays use the induction principle discovered by Galileo Ferraris in the late 19th century. The magnetic system in induction disc overcurrent relays is designed to detect overcurrents in a power system and operate with a pre-determined time delay when certain overcurrent limits have been reached. In order to operate, the magnetic system in the relays produces torque that acts on a metal disc to make contact, according to the following basic current/torque equation:
Where and are the two fluxes and is the phase angle between the fluxes
The following important conclusions can be drawn from the above equation.
Two alternating fluxes with a phase shift are needed for torque production.
Maximum torque is produced when the two alternating fluxes are 90 degrees apart.
The resultant torque is steady and not a function of time.
The relay's primary winding is supplied from the power systems current transformer via a plug bridge, which is called the plug setting multiplier (psm). Usually seven equally spaced tappings or operating bands determine the relays sensitivity. The primary winding is located on the upper electromagnet. The secondary winding has connections on the upper electromagnet that are energised from the primary winding and connected to the lower electromagnet. Once the upper and lower electromagnets are energised they produce eddy currents that are induced onto the metal disc and flow through the flux paths. This relationship of eddy currents and fluxes creates torque proportional to the input current of the primary winding, due to the two flux paths being out of phase by 90°.
In an overcurrent condition, a value of current will be reached that overcomes the control spring pressure on the spindle and the braking magnet, causing the metal disc to rotate towards the fixed contact. This initial movement of the disc is also held off to a critical positive value of current by small slots that are often cut into the side of the disc. The time taken for rotation to make the contacts is not only dependent on current but also the spindle backstop position, known as the time multiplier (tm). The time multiplier is divided into 10 linear divisions of the full rotation time.
Providing the relay is free from dirt, the metal disc and the spindle with its contact will reach the fixed contact, thus sending a signal to trip and isolate the circuit, within its designed time and current specifications. Drop off current of the relay is much lower than its operating value, and once reached the relay will be reset in a reverse motion by the pressure of the control spring governed by the braking magnet.
Static
Application of electronic amplifiers to protective relays was described as early as 1928, using vacuum tube amplifiers and continued up to 1956. Devices using electron tubes were studied but never applied as commercial products, because of the limitations of vacuum tube amplifiers. A relatively large standby current is required to maintain the tube filament temperature; inconvenient high voltages are required for the circuits, and vacuum tube amplifiers had difficulty with incorrect operation due to noise disturbances.
Static relays have no or few moving parts, and became practical with the introduction of the transistor. Measuring elements of static relays have been successfully and economically built up from diodes, zener diodes, avalanche diodes, unijunction transistors, p-n-p and n-p-n bipolar transistors, field effect transistors or their combinations. Static relays offer the advantage of higher sensitivity than purely electromechanical relays, because power to operate output contacts is derived from a separate supply, not from the signal circuits. Static relays eliminated or reduced contact bounce, and could provide fast operation, long life and low maintenance.
Digital
Digital protective relays were in their infancy during the late 1960s. An experimental digital protection system was tested in the lab and in the field in the early 1970s. Unlike the relays mentioned above, digital protective relays have two main parts: hardware and software. The world's first commercially available digital protective relay was introduced to the power industry in 1984 by Schweitzer Engineering Laboratories (SEL) based in Pullman, Washington. In spite of the developments of complex algorithms for implementing protection functions the microprocessor based-relays marketed in the 1980s did not incorporate them.
A microprocessor-based digital protection relay can replace the functions of many discrete electromechanical instruments. These relays convert voltage and currents to digital form and process the resulting measurements using a microprocessor. The digital relay can emulate functions of many discrete electromechanical relays in one device, simplifying protection design and maintenance. Each digital relay can run self-test routines to confirm its readiness and alarm if a fault is detected. Digital relays can also provide functions such as communications (SCADA) interface, monitoring of contact inputs, metering, waveform analysis, and other useful features. Digital relays can, for example, store multiple sets of protection parameters, which allows the behavior of the relay to be changed during maintenance of attached equipment. Digital relays also can provide protection strategies impossible to implement with electromechanical relays. This is particularly so in long-distance high voltage or multi-terminal circuits or in lines that are series or shunt compensated They also offer benefits in self-testing and communication to supervisory control systems.
Numerical
The distinction between digital and numerical protection relay rests on points of fine technical detail, and is rarely found in areas other than Protection. Numerical relays are the product of the advances in technology from digital relays. Generally, there are several different types of numerical protection relays. Each type, however, shares a similar architecture, thus enabling designers to build an entire system solution that is based on a relatively small number of flexible components. They use high speed processors executing appropriate algorithms. Most numerical relays are also multifunctional and have multiple setting groups each often with tens or hundreds of settings.
Relays by functions
The various protective functions available on a given relay are denoted by standard ANSI device numbers. For example, a relay including function 51 would be a timed overcurrent protective relay.
Overcurrent relay
An overcurrent relay is a type of protective relay which operates when the load current exceeds a pickup value. It is of two types: instantaneous over current (IOC) relay and definite time overcurrent (DTOC) relay.
The ANSI device number is 50 for an IOC relay or a DTOC relay. In a typical application, the over current relay is connected to a current transformer and calibrated to operate at or above a specific current level. When the relay operates, one or more contacts will operate and energize to trip a circuit breaker. The DTOC relay has been used extensively in the United Kingdom but its inherent issue of operating slower for faults closer to the source led to the development of the IDMT relay.
Definite time over-current relay
A definite time over-current (DTOC) relay is a relay that operates after a definite period of time once the current exceeds the pickup value. Hence, this relay has current setting range as well as time setting range.
Instantaneous over-current relay
An instantaneous over-current relay is an overcurrent relay which has no intentional time delay for operation. The contacts of the relay are closed instantly when the current inside the relay rises beyond the operational value. The time interval between the instant pick-up value and the closing contacts of the relay is very low. It has low operating time and starts operating instantly when the value of current is more than the relay setting. This relay operates only when the impedance between the source and the relay is less than that provided in the section.
Inverse-time over-current relay
An inverse-time over-current (ITOC) relay is an overcurrent relay which operates only when the magnitude of their operating current is inversely proportional to the magnitude of the energize quantities. The operating time of relay decreases with the increases in the current. The operation of the relay depends on the magnitude of the current.
Inverse definite minimum time relay
The inverse definite minimum time (IDMT) relay are protective relays which were developed to overcome the shortcomings of the definite time overcurrent (DTOC) relays.
If the source impedance remains constant and the fault current changes appreciably as we move away from the relay then it is advantageous to use IDMT overcurrent protection to achieve high speed protection over a large section of the protected circuit. However, if the source impedance is significantly larger than the feeder impedance then the characteristic of the IDMT relay cannot be exploited and DTOC may be utilized. Secondly if the source impedance varies and becomes weaker with less generation during light loads then this leads to slower clearance time hence negating the purpose of the IDMT relay.
IEC standard 60255-151 specifies the IDMT relay curves as shown below. The four curves in Table 1 are derived from the now withdrawn British Standard BS 142. The other five, in Table 2, are derived from the ANSI standard C37.112.
While it is more common to use IDMT relays for current protection it is possible to utilize IDMT mode of operation for voltage protection. It is possible to program customised curves in some protective relays and other manufacturers have special curves specific to their relays. Some numerical relays can be used to provide inverse time overvoltage protection or negative sequence overcurrent protection.
Ir = is the ratio of the fault current to the relay setting current or a Plug Setting Multiplier. "Plug" is a reference from the electromechanical relay era and were available in discrete steps. TD is the Time Dial setting.
The above equations result in a "family" of curves as a result of using different time multiplier setting (TMS) settings. It is evident from the relay characteristic equations that a larger TMS will result in a slower clearance time for a given PMS (I) value.
Distance relay
Distance relays, also known as impedance relay, differ in principle from other forms of protection in that their performance is not governed by the magnitude of the current or voltage in the protected circuit but rather on the ratio of these two quantities. Distance relays are actually double actuating quantity relays with one coil energized by voltage and other coil by current. The current element produces a positive or pick up torque while the voltage element produces a negative or reset torque. The relay operates only when the V/I ratio falls below a predetermined value (or set value). During a fault on the transmission line the fault current increases and the voltage at the fault point decreases. The V/I
ratio is measured at the location of CTs and PTs. The voltage at the PT location depends on the distance between the PT and the fault. If the measured voltage is lesser, that means the fault is nearer and vice versa. Hence the protection called Distance relay. The load flowing through the line appears as an impedance to the relay and sufficiently large loads (as impedance is inversely proportional to the load) can lead to a trip of the relay even in the absence of a fault.
Current differential protection scheme
A differential scheme acts on the difference between current entering a protected zone (which may be a bus bar, generator, transformer or other apparatus) and the current leaving that zone. A fault outside the zone gives the same fault current at the entry and exit of the zone, but faults within the zone show up as a difference in current.
"The differential protection is 100% selective and therefore only responds to faults within its protected zone. The boundary of the protected zone is uniquely defined by the location of the current transformers. Time grading with other protection systems is therefore not required, allowing for tripping without additional delay. Differential protection is therefore suited as fast main protection for all important plant items."
Differential protection can be used to provide protection for zones with multiple terminals and can be used to protect lines, generators, motors, transformers, and other electrical plant.
Current transformers in a differential scheme must be chosen to have near-identical response to high overcurrents. If a "through fault" results in one set of current transformers saturating before another, the zone differential protection will see a false "operate" current and may false trip.
GFCI (ground fault circuit interrupter) circuit breakers combine overcurrent protection and differential protection (non-adjustable) in standard, commonly available modules.
Directional relay
A directional relay uses an additional polarizing source of voltage or current to determine the direction of a fault. Directional elements respond to the phase shift between a polarizing quantity and an operate quantity. The fault can be located upstream or downstream of the relay's location, allowing appropriate protective devices to be operated inside or outside of the zone of protection.
Synchronism check
A synchronism checking relay provides a contact closure when the frequency and phase of two sources are similar to within some tolerance margin. A "synch check" relay is often applied where two power systems are interconnected, such as at a switchyard connecting two power grids, or at a generator circuit breaker to ensure the generator is synchronized to the system before connecting it.
Power source
The relays can also be classified on the type of power source that they use to work.
Self-powered relays operate on energy derived from the protected circuit, through the current transformers used to measure line current, for example. This eliminates the cost and reliability question of a separate supply.
Auxiliary powered relays rely on a battery or external ac supply. Some relays can use either AC or DC. The auxiliary supply must be highly reliable during a system fault.
Dual powered relays can be also auxiliary powered, so all batteries, chargers and other external elements are made redundant and used as a backup.
| Technology | Electrical protective devices | null |
27072071 | https://en.wikipedia.org/wiki/Unit%20circle | Unit circle | In mathematics, a unit circle is a circle of unit radius—that is, a radius of 1. Frequently, especially in trigonometry, the unit circle is the circle of radius 1 centered at the origin (0, 0) in the Cartesian coordinate system in the Euclidean plane. In topology, it is often denoted as because it is a one-dimensional unit -sphere.
If is a point on the unit circle's circumference, then and are the lengths of the legs of a right triangle whose hypotenuse has length 1. Thus, by the Pythagorean theorem, and satisfy the equation
Since for all , and since the reflection of any point on the unit circle about the - or -axis is also on the unit circle, the above equation holds for all points on the unit circle, not only those in the first quadrant.
The interior of the unit circle is called the open unit disk, while the interior of the unit circle combined with the unit circle itself is called the closed unit disk.
One may also use other notions of "distance" to define other "unit circles", such as the Riemannian circle; see the article on mathematical norms for additional examples.
In the complex plane
In the complex plane, numbers of unit magnitude are called the unit complex numbers. This is the set of complex numbers such that When broken into real and imaginary components this condition is
The complex unit circle can be parametrized by angle measure from the positive real axis using the complex exponential function, (See Euler's formula.)
Under the complex multiplication operation, the unit complex numbers form a group called the circle group, usually denoted In quantum mechanics, a unit complex number is called a phase factor.
Trigonometric functions on the unit circle
The trigonometric functions cosine and sine of angle may be defined on the unit circle as follows: If is a point on the unit circle, and if the ray from the origin to makes an angle from the positive -axis, (where counterclockwise turning is positive), then
The equation gives the relation
The unit circle also demonstrates that sine and cosine are periodic functions, with the identities
for any integer .
Triangles constructed on the unit circle can also be used to illustrate the periodicity of the trigonometric functions. First, construct a radius from the origin to a point on the unit circle such that an angle with is formed with the positive arm of the -axis. Now consider a point and line segments . The result is a right triangle with . Because has length , length , and has length 1 as a radius on the unit circle, and . Having established these equivalences, take another radius from the origin to a point on the circle such that the same angle is formed with the negative arm of the -axis. Now consider a point and line segments . The result is a right triangle with . It can hence be seen that, because , is at in the same way that P is at . The conclusion is that, since is the same as and is the same as , it is true that and . It may be inferred in a similar manner that , since and . A simple demonstration of the above can be seen in the equality .
When working with right triangles, sine, cosine, and other trigonometric functions only make sense for angle measures more than zero and less than . However, when defined with the unit circle, these functions produce meaningful values for any real-valued angle measure – even those greater than 2. In fact, all six standard trigonometric functions – sine, cosine, tangent, cotangent, secant, and cosecant, as well as archaic functions like versine and exsecant – can be defined geometrically in terms of a unit circle, as shown at right.
Using the unit circle, the values of any trigonometric function for many angles other than those labeled can be easily calculated by hand using the angle sum and difference formulas.
Complex dynamics
The Julia set of discrete nonlinear dynamical system with evolution function:
is a unit circle. It is a simplest case so it is widely used in the study of dynamical systems.
| Mathematics | Trigonometry | null |
27077017 | https://en.wikipedia.org/wiki/Flocculent%20spiral%20galaxy | Flocculent spiral galaxy | A flocculent spiral galaxy is a type of spiral galaxy. Unlike the well-defined spiral architecture of a grand design spiral galaxy, flocculent (meaning "flaky") galaxies are patchy, with discontinuous spiral arms. Self-propagating star formation is the apparent explanation for the structure of flocculent spirals. Approximately 30% of spirals are flocculent, 10% are grand design, and the rest are referred to as "multi-armed". The multiple-arm type is sometimes grouped into the flocculent category.
The prototypical flocculent spiral is NGC 2841.
List of flocculent spiral galaxies
| Physical sciences | Galaxy classification | Astronomy |
1972407 | https://en.wikipedia.org/wiki/Neritic%20zone | Neritic zone | The neritic zone (or sublittoral zone) is the relatively shallow part of the ocean above the drop-off of the continental shelf, approximately in depth.
From the point of view of marine biology it forms a relatively stable and well-illuminated environment for marine life, from plankton up to large fish and corals, while physical oceanography sees it as where the oceanic system interacts with the coast.
Definition (marine biology), context, extra terminology
In marine biology, the neritic zone, also called coastal waters, the coastal ocean or the sublittoral zone, refers to that zone of the ocean where sunlight reaches the ocean floor, that is, where the water is never so deep as to take it out of the photic zone.
It extends from the low tide mark to the edge of the continental shelf, with a relatively shallow depth extending to about 200 meters (660 feet).
Above the neritic zone lie the intertidal (or eulittoral) and supralittoral zones; below it the continental slope begins, descending from the continental shelf to the abyssal plain and the pelagic zone.
Within the neritic, marine biologists also identify the following:
The infralittoral zone is the algal-dominated zone down to around five metres below the low water mark.
The circalittoral zone is the region beyond the infralittoral, which is dominated by sessile animals such as oysters.
The subtidal zone is the region of the neritic zone which is below the intertidal zone, therefore never exposed to the atmosphere.
Physical characteristics
The neritic zone is covered with generally well-oxygenated water, receives plenty of sunlight, is relatively stable temperature, has low water pressure and stable salinity levels, making it highly suitable for photosynthetic life.
There are several different areas or zones in the ocean. The area along the bottom of any body of water from the shore to the deepest abyss is called the benthic zone. It is where decomposed organic debris (also known as ocean 'snow') has settled to form a sediment layer. All photosynthetic life needs light to grow and how far out into the ocean light can still penetrate through the water column to the floor or benthic zone is what defines the neritic zone. That photic zone, or area where light can penetrate through the water column, is usually above ~100 meters (~328 feet). Some coastal areas have a long area of shallow water that extends far out beyond the landmass into the water and others, for example islands that have formed from ancient volcanos where the 'shelf' or edge of the land mass is very steep, have a very short neritic zone.
Life forms
The above characteristics make the neritic zone the location of the majority of sea life.
The result is high primary production by photosynthetic life such as phytoplankton and floating sargassum;
zooplankton, free-floating creatures ranging from microscopic foraminiferans to small fish and shrimp, feed on the phytoplankton (and one another);
both trophic levels in turn form the base of the food chain (or, more properly, web) that supports most of the world's great wild fisheries.
Corals are also mostly found in the neritic zone, where they are more common than in the intertidal zone as they have less change to deal with.
Definition (physical oceanography)
In physical oceanography, the sublittoral zone refers to coastal regions with significant tidal flows and energy dissipation, including non-linear flows, internal waves, river outflows and ocean fronts.
As in marine biology, this zone typically extends to the edge of the continental shelf.
| Physical sciences | Oceanography | Earth science |
1973155 | https://en.wikipedia.org/wiki/Parasaurolophus | Parasaurolophus | Parasaurolophus (; meaning "beside crested lizard" in reference to Saurolophus) is a genus of hadrosaurid "duck-billed" dinosaur that lived in what is now western North America and possibly Asia during the Late Cretaceous period, about 76.9–73.5 million years ago. It was a large herbivore that could reach over long and weigh over , and were able to move as a biped and a quadruped. Three species are universally recognized: P. walkeri (the type species), P. tubicen, and the short-crested P. cyrtocristatus. Additionally, a fourth species, P. jiayinensis, has been proposed, although it is more commonly placed in the separate genus Charonosaurus. Remains are known from Alberta, New Mexico, and Utah, as well as possibly Heilongjiang if Charonosaurus is in fact part of the genus. The genus was first described in 1922 by William Parks from a skull and partial skeleton found in Alberta.
Parasaurolophus was a hadrosaurid, part of a diverse family of large Late Cretaceous ornithopods that are known for their range of bizarre head adornments, which were likely used for communication and increased hearing. This genus is known for its large, elaborate cranial crest, which forms a long curved tube projecting upwards and back from the skull in its largest form. Charonosaurus from China, which may have been its closest relative, had a similar skull and a potentially similar crest. Visual recognition of both species and sex, acoustic resonance, and thermoregulation have been proposed as functional explanations for the crest. It is one of the rarer hadrosaurids, known from only a handful of good specimens.
Discovery and naming
Meaning "near crested lizard", the name Parasaurolophus is derived from the Greek words para/παρα ("beside" or "near"), saurus/ ("lizard"), and lophos/λοφος ("crest"). It is based on ROM 768, a skull and partial skeleton missing most of the tail and the back legs below the knees, which was found by a field party from the University of Toronto in 1920 near Sand Creek along the Red Deer River in Alberta. These rocks are now known as the Campanian age Late Cretaceous Dinosaur Park Formation. William Parks named the specimen P. walkeri in honor of Sir Byron Edmund Walker, the chairman of the Board of Trustees of the Royal Ontario Museum. Parasaurolophus remains are rare in Alberta, with only one other partial skull that is possibly from the Dinosaur Park Formation and three Dinosaur Park specimens lacking their skulls that possibly belong to the genus. In some faunal lists, there is a mention of possible P. walkeri material in the Hell Creek Formation of Montana, a rock unit of the late Maastrichtian age. This occurrence is not noted by Sullivan and Williamson in their 1999 review of the genus and has not been further elaborated upon elsewhere.
In 1921, Charles H. Sternberg recovered a partial skull (PMU.R1250) from what is now known as the slightly younger Kirtland Formation in San Juan County, New Mexico. This specimen was sent to Uppsala, where Carl Wiman described it as a second species, P. tubicen, in 1931. The specific epithet is derived from the Latin word tǔbǐcěn, meaning "trumpeter". A second, nearly complete P. tubicen skull (NMMNH P-25100) was found in New Mexico in 1995. Using computed tomography of this skull, Robert Sullivan and Thomas Williamson gave the genus a monographic treatment in 1999 that covered aspects of its anatomy and taxonomy, as well as the functions of its crest. Williamson later published an independent review of the remains that disagreed with the taxonomic conclusions.
John Ostrom described another good specimen (FMNH P27393) from New Mexico as P. cyrtocristatus in 1961. It includes a partial skull with a short, rounded crest and much of the postcranial skeleton except for the feet, neck, and parts of the tail. Its specific name is derived from the Latin words curtus, meaning "shortened" and cristatus, meaning "crested". The specimen was found in either the top of the Fruitland Formation or, more likely, the base of the overlying Kirtland Formation. The range of this species was described in 1979, when David B. Weishampel and James A. Jensen described a partial skull with a similar crest (BYU 2467) from the Campanian age Kaiparowits Formation of Garfield County, Utah. Since then, another skull has been found in Utah with the short, rounded P. cyrtocristatus crest morphology.
Species
Parasaurolophus is known from three certain species: P. walkeri, P. tubicen, and P. cyrtocristatus. All of them can be clearly distinguished from each other and have many differences. The first named species, therefore the type, is P. walkeri. One certain specimen from the Dinosaur Park Formation is referred to it, but many more are almost certainly referable. Like stated above, it is different from the other two species, with it having a simpler internal structure than P. tubicen, along with a straighter crest and different internal structuring than P. cyrtocristatus.
The next named species is P. tubicen, which is the largest of the Parasaurolophus species. It lived in New Mexico, where three specimens are known, and can be differentiated from its other species. It possesses a long and straight crest, with a very complex interior compared to the other species. All known specimens of P. tubicen come from the De-Na-Zin Member of the Kirtland Formation.
In 1961, the third species, P. cyrtocristatus was named by John Ostrom. Its three known specimens have been found in the Fruitland and Kaiparowits formations of Utah and New Mexico. The second specimen, the first known from the Kaiparowits Formation, was originally unassigned to a specific taxon. Of the Parasaurolophus species, P. cyrtocristatus is the smallest and has the most curved crest. Because of its possession of the two above features, it has often been speculated that it was a female of P. walkeri or P. tubicen, which were all thought to be males, although P. tubicen lived approximately a million years later. As noted by Thomas Williamson, the type material of P. cyrtocristatus is about 72% the size of P. tubicen, close to the size at which other lambeosaurines are interpreted to begin showing definitive sexual dimorphism in their crests (~70% of adult size). Even though many scientists have supported the possible fact of P. cyrtocristatus being a female, many other studies have found that it is not because of the differences in age, distribution, and the large differences in the crest and its internal structure.
A study published in PLoS ONE in 2014 found that one more species could be referred to Parasaurolophus. This study, led by Xing, found Charonosaurus jiayensis was actually nested deeply inside Parasaurolophus, which created the new species P. jiayensis. If this species is indeed inside Parasaurolophus, then the genus therefore lasted until the K-Pg extinction and is known from two continents.
Description
Like most dinosaurs, the skeleton of Parasaurolophus is incompletely known. The length of the type specimen of P. walkeri is estimated at , and allometry-based body mass estimates indicate that a long individual would have weighed more than . Gregory S. Paul estimated that an average adult individual of the type species would measure long and weigh .
Its skull is about long, including the crest, whereas the type skull of P. tubicen is over long, indicating it was a larger animal. Its single known arm was relatively short for a hadrosaurid, with a short but wide shoulder blade. The thighbone measures long in P. walkeri and is robust for its length when compared to other hadrosaurids. The upper arm and pelvic bones were also heavily built.
Like other hadrosaurids, it was able to walk on either two legs or four. It probably preferred to forage for food on four legs, but ran on two. The neural spines of the vertebrae were tall, as was common in lambeosaurines. At their tallest over the hips, they increased the height of the back. Skin impressions are known for P. walkeri, showing uniform tubercle-like scales, but no larger structures.
Skull
The most noticeable feature was the cranial crest that protruded from the rear of the head and was made up of the premaxilla and nasal bones. The crest was hollow, with distinct tubes leading from each nostril to the end of the crest before reversing direction and heading back down the crest and into the skull. The tubes were simplest in P. walkeri, and more complex in P. tubicen, where some tubes were blind and others met and separated. While P. walkeri and P. tubicen had long crests with slight curvature, P. cyrtocristatus had a short crest with a more circular profile.
Classification
As its name implies, Parasaurolophus was initially thought to be closely related to Saurolophus because of its superficially similar crest. However, it was soon reassessed as a member of the lambeosaurine subfamily of hadrosaurids—Saurolophus is a hadrosaurine. It is usually interpreted as a separate offshoot of the lambeosaurines, distinct from the helmet-crested Corythosaurus, Hypacrosaurus, and Lambeosaurus. Its closest known relative appears to be Charonosaurus, a lambeosaurine with a similar skull (but no complete crest yet) from the Amur region of northeastern China. The two may form the clade Parasaurolophini. P. cyrtocristatus, with its short, rounded crest, may be the most basal of the three known Parasaurolophus species or it may represent subadult or female specimens of P. tubicen.
The following cladogram is after the 2007 redescription of Lambeosaurus magnicristatus (Evans and Reisz, 2007):
Paleobiology
Diet and feeding
As a hadrosaurid, Parasaurolophus was a large bipedal and quadrupedal herbivore, eating plants with a sophisticated skull that permitted a grinding motion analogous to chewing. Its teeth were continually being replaced and were packed into dental batteries containing hundreds of teeth, but only a relative handful of which were in use at any time. It used its beak to crop plant material, which was held in the jaws by a cheek-like organ. Vegetation could have been taken from the ground up to a height of around . As noted by Robert Bakker, lambeosaurines have narrower beaks than hadrosaurines, implying that Parasaurolophus and its relatives could feed more selectively than their broad-beaked, crestless counterparts. Parasaurolophus had a diet consisting of leaves, twigs, and pine needles which would imply that it was a browser.
Growth
Parasaurolophus is known from many adult specimens, and a juvenile described in 2013, numbered RAM 140000 and nicknamed Joe, after a volunteer at the Raymond M. Alf Museum of Paleontology (RAM). The juvenile was discovered in the Kaiparowits Formation in 2009. Excavated by the joint expedition by museum and The Webb Schools, the juvenile has been identified as around only one year old when it died. Referred to Parasaurolophus sp., the juvenile is the most complete, as well as youngest Parasaurolophus ever found, and measures . This individual fits neatly into the currently known Parasaurolophus growth stages, and lived approximately 75 million years ago. Even though no complete skull of the intermediate age between RAM 14000 and adult Parasaurolophus has been found yet, a partial braincase of about the right size is known. At 25% of the total adult size, the juvenile show that crest growth of Parasaurolophus began sooner than in related genera, such as Corythosaurus. It has been suggested that Parasaurolophus adults bore such large crests, especially when compared to the related Corythosaurus, because of this difference in age between when their crests started developing. The crest of the juvenile is not long and tubular like the adults, but low and hemispherical.
The skull of RAM 14000 is almost complete, with the left side only lacking a piece of the maxilla. However, the skull was split down the middle by erosion, possibly when it was resting on the bottom of a river bed. The two sides are displaced slightly, with some bones of the right being moved off the main block, also by erosion. After reconstruction, the skull viewed from the side resembles other juvenile lambeosaurines found, being roughly a trapezoid in shape.
A partial cranial endocast for RAM 14000 was reconstructed from CT scan data, the first ever for a Parasaurolophus of any ontogenetic stage. The endocast was reconstructed in two sections, one on the portion of the braincase articulated with the left half of the skull and the remainder on the disarticulated portion of the braincase. Their relative position was then approximated based on cranial landmarks and comparison with other hadrosaurids. Because of weathering, many of the smaller neural canals and foramina could not be identified for certain.
Cranial crest
Many hypotheses have been advanced as to what functions the cranial crest of Parasaurolophus performed, but most have been discredited. It is now believed that it may have had several functions: visual display for identifying species and sex, sound amplification for communication, and thermoregulation. It is not clear which was most significant at what times in the evolution of the crest and its internal nasal passages.
Differences in crests
As for other lambeosaurines, it is believed that the cranial crest of Parasaurolophus changed with age and was a sexually dimorphic characteristic in adults. James Hopson, one of the first researchers to describe lambeosaurine crests in terms of such distinctions, suggested that P. cyrtocristatus, with its small crest, was the female form of P. tubicen. Thomas Williamson suggested it was the juvenile form. Neither hypothesis became widely accepted. As only six good skulls, one juvenile braincase, and one recently discovered juvenile skull are known, additional material will help clear up these potential relationships. Williamson noted that in any case, juvenile Parasaurolophus probably had small, rounded crests like P. cyrtocristatus, that probably grew faster as individuals approached sexual maturity. Recent restudy of a juvenile braincase previously assigned to Lambeosaurus, now assigned to Parasaurolophus, provides evidence that a small tubular crest was present in juveniles. This specimen preserves a small upward flaring of the frontal bones that was similar to but smaller than what is seen in adult specimens; in adults, the frontals formed a platform that supported the base of the crest. This specimen also indicates that the growth of the crest in Parasaurolophus and the facial profile of juvenile individuals differed from the Corythosaurus-Hypacrosaurus-Lambeosaurus model, in part because the crest of Parasaurolophus lacks the thin bony 'coxcomb' that makes up the upper portion of the crest of the other three lambeosaurines.
Rejected function hypotheses
Many early suggestions focused on adaptations for an aquatic lifestyle, following the hypothesis that hadrosaurids were amphibious, a common line of thought until the 1960s. Thus, Alfred Sherwood Romer proposed it served as a snorkel, Martin Wilfarth that it was an attachment for a mobile proboscis used as a breathing tube or for food gathering, Charles M. Sternberg that it served as an airtrap to keep water out of the lungs, and Ned Colbert that it served as an air reservoir for prolonged stays underwater.
Other proposals were more mechanical in nature. William Parks, in 1922, suggested that the crest was joined to the vertebrae above the shoulders by ligaments or muscles, and helped with moving and supporting the head. This is unlikely, because in all modern archosaurs, the nuchal ligament attaches to the neck or base of the skull. Othenio Abel proposed it was used as a weapon in combat among members of the same species, and Andrew Milner suggested that it could be used as a foliage deflector, like the helmet crest (called a 'casque') of the cassowary. Still, other proposals made housing specialized organs the major function. Halszka Osmólska suggested that it housed salt glands, and John Ostrom suggested that it housed expanded areas for olfactory tissue and much improved sense of smell of the lambeosaurines, which had no obvious defensive capabilities.
Most of these hypotheses have been discredited or rejected. For example, there is no hole at the end of the crest for a snorkeling function. There are no muscle scars for a proboscis and it is dubious that an animal with a beak would need one. As a proposed airlock, it would not have kept out water. The proposed air reservoir would have been insufficient for an animal the size of Parasaurolophus. Other hadrosaurids had large heads without needing large hollow crests to serve as attachment points for supporting ligaments. Also, none of the proposals explain why the crest has such a shape, why other lambeosaurines should have crests that look much different but perform a similar function, how crestless or solid-crested hadrosaurids got along without such capabilities, or why some hadrosaurids had solid crests. These considerations particularly impact hypotheses based on increasing the capabilities of systems already present in the animal, such as the salt gland and olfaction hypotheses, and indicate that these were not primary functions of the crest. Additionally, work on the nasal cavity of lambeosaurines shows that olfactory nerves and corresponding sensory tissue were largely outside the portion of the nasal passages in the crest, so the expansion of the crest had little to do with the sense of smell.
Temperature regulation hypothesis
The large surface area and vascularization of the crest also suggests a thermoregulatory function. The first to propose the cranial crests of lambeosaurines related to temperature regulation was Wheeler (1978). He proposed that there was a nerve connection between the crest and the brain, so that the latter could be cooled by the former. The next people to publish a related idea were Teresa Maryańska and Osmólska, who realized that like modern lizards, dinosaurs could have possessed salt glands, and cooled off by osmo-regulation. In 2006 Evans published an argument about the functions of lambeosaurine crests, and supported why this could be a causing factor for the evolution of the crest.
Behavioral hypotheses
Parasaurolophus is often hypothesized to have used its crest as a resonating chamber to produce low frequency sounds to alert other members of a group or its species. This function was originally suggested by Wiman in 1931 when he described P. tubicen. He noted that the crest's internal structures are similar to those of a swan and theorized that an animal could use its elongated nasal passages to create noise. However, the nasal tubes of Hypacrosaurus, Corythosaurus, and Lambeosaurus are much more variable and complicated than the airway of Parasaurolophus. A large amount of material and data supports the hypothesis that the large, tubular crest of Parasaurolophus was a resonating chamber. Weishampel in 1981 suggested that Parasaurolophus made noises ranging between the frequencies 55 and 720 Hz, although there was some difference in the range of individual species because of the crest size, shape, and nasal passage length, most obvious in P. cyrtocristatus (interpreted as a possible female). Hopson found that there is anatomical evidence that hadrosaurids had a strong hearing. There is at least one example, in the related Corythosaurus, of a slender stapes (reptilian ear bone) in place, which combined with a large space for an eardrum implies a sensitive middle ear. Furthermore, the hadrosaurid lagena is elongate like a crocodilian's, indicating that the auditory portion of the inner ear was well-developed. Based on the similarity of hadrosaurid inner ears to those of crocodiles, he also proposed that adult hadrosaurids were sensitive to high frequencies, such as their offspring might produce. According to Weishampel, this is consistent with parents and offspring communicating.
Computer modeling of a well-preserved specimen of P. tubicen, with more complex air passages than those of P. walkeri, has allowed the reconstruction of the possible sound its crest produced. The main path resonates at around 30 Hz, but the complicated sinus anatomy causes peaks and valleys in the sound.
The other main behavioral theory is that the crest was used for intra-species recognition. This means that the crest could have been used for species recognition, as a warning signal, and for other, non-sexual uses. These could have been some of the reasons crests evolved in Parasaurolophus and other hadrosaurids. Instead, social and physiological functions have become more supported as function(s) of the crest, focusing on visual and auditory identification and communication. As a large object, the crest has clear value as a visual signal and sets this animal apart from its contemporaries. The large size of hadrosaurid eye sockets and the presence of sclerotic rings in the eyes imply acute vision and diurnal habits, evidence that sight was important to these animals. If, as is commonly illustrated, a skin frill extended from the crest to the neck or back, the proposed visual display would have been even showier. As is suggested by other lambeosaurine skulls, the crest of Parasaurolophus likely permitted both species identification (such as separating it from Corythosaurus or Lambeosaurus) and sexual identification by shape and size.
Soft tissue frill
Barnum Brown (1912) noted the presence of fine striations near the back of the crest that he hypothesized could be associated with the presence of a frill of skin, comparable to the one found in the modern basilisk lizard. His hypothesis was seemingly supported by skin preserved above the neck and back of Corythosaurus and Edmontosaurus. Subsequently, reconstructions of Parasaurolophus with a substantial frill of skin between the crest and neck appeared in influential paleoart including murals by Charles R. Knight and in the Walt Disney animated film, Fantasia. This led to the frill being depicted in many other sources, though the advent of the now-debunked "snorkel" hypothesis, and conflation of the frill hypothesis with the idea that the crest serves as an anchor point for neck ligaments, along with lack of strong evidence for its presence, has seen it fall out of favor in most modern depictions.
Paleopathology
P. walkeri is known from one specimen which might contain a pathology. The skeleton shows a v-shaped gap or notch in the vertebrae at the base of the neck. Originally thought to be pathologic, Parks published a second interpretation of this, as a ligament attachment to support the head. The crest would attach to the gap via muscles or ligaments, and be used to support the head while bearing a frill, like predicted to exist in some hadrosaurids. One other possibility, is that during preparation, the specimen was damaged, creating the possible pathology. The notch, however, is still considered more likely to be a pathology, even though some illustrations of Parasaurolophus restore the skin flap.
Another possible pathology was noticed by Parks, and from around the notch. In the fourth, fifth, and sixth vertebrae, directly anterior to the notch, the neural spines were damaged. The fourth had an obvious fracture, with the other two possessing a swelling at the base of the break.
Analysis of the pathology undertaken by Bertozzo et al., published in December 2020, suggests the pathology to the shoulder and thoracic ribs in the holotype of P. walkeri was plausibly the result of the dinosaur being hit by a falling tree, perhaps during a severe storm. Based on the regrowth of bone, it is suggested that the hadrosaur survived for at least one to four months to perhaps years after being injured. None of the pathologies on the holotype individual are believed to have caused or contributed to its death.
Paleoecology
Alberta
Parasaurolophus walkeri, from the Dinosaur Park Formation, was a member of a diverse and well-documented fauna of prehistoric animals, including well-known dinosaurs such as the horned Centrosaurus, Chasmosaurus, and Styracosaurus; ornithomimids Struthiomimus; fellow duckbills Gryposaurus and Corythosaurus; tyrannosaurids Gorgosaurus and Daspletosaurus; and armored Edmontonia, Euoplocephalus and Dyoplosaurus. It was a rare constituent of this fauna. The Dinosaur Park Formation is interpreted as a low-relief setting of rivers and floodplains that became more swampy and influenced by marine conditions over time as the Western Interior Seaway transgressed westward. The climate was warmer than present-day Alberta, without frost, but with wetter and drier seasons. Conifers were apparently the dominant canopy plants, with an understory of ferns, tree ferns, and angiosperms.
Some of the less common hadrosaurs in the Dinosaur Park Formation of Dinosaur Provincial Park, such as Parasaurolophus, may represent the remains of individuals who died while migrating through the region. They might also have had a more upland habitat where they may have nested or fed. The presence of Parasaurolophus and Kritosaurus in northern latitude fossil sites may represent faunal exchange between otherwise distinct northern and southern biomes in Late Cretaceous North America. Both taxa are uncommon outside of the southern biome, where, along with Pentaceratops, they are predominate members of the fauna.
New Mexico
In the Fruitland Formation of New Mexico, P. cyrtocristatus shared its habitat with other ornithischians and theropods. Specifically, its contemporaries were the ceratopsian Pentaceratops sternbergii; the pachycephalosaur Stegoceras novomexicanum; and some unidentified fossils belonging to Tyrannosauridae, ?Ornithomimus, ?Troodontidae, ?Saurornitholestes langstoni, ?Struthiomimus, Ornithopoda, ?Chasmosaurus, ?Corythosaurus, Hadrosaurinae, Hadrosauridae, and Ceratopsidae. When Parasaurolophus existed, the Fruitland Formation was swampy, positioned in the lowlands, and close to the shore of the Cretaceous Interior Seaway. The lowermost part of the Fruitland Formation is just younger than 75.56 ± 0.41 mya, with the uppermost boundary dating to 74.55 ± 0.22 mya.
Existing slightly later than the species from the Fruitland Formation, P. tubicen is also found in New Mexico, in the Kirtland Formation. Numerous vertebrate groups are from this formation, including fishes, crurotarsans, ornithischians, saurischians, pterosaurs, and turtles. The fishes are represented by the two species Melvius chauliodous and Myledalphus bipartitus. The crurotarsans include Brachychampsa montana and Denazinosuchus kirtlandicus. Ornithischians from the formation are represented by the hadrosaurids Anasazisaurus horneri, Naashoibitosaurus ostromi, Kritosaurus navajovius, and P. tubicen; the ankylosaurids Ahshislepelta minor and Nodocephalosaurus kirtlandensis; the ceratopsians Pentaceratops sternbergii and Titanoceratops ouranos; and the pachycephalosaurs Stegoceras novomexicanum and Sphaerotholus goodwini. Saurischians include the tyrannosaurid Bistahieversor sealeyi; the ornithomimid Ornithomimus sp.; and the troodontid "Saurornitholestes" robustus. One pterosaur is known, named Navajodactylus boerei. Turtles are fairly plentiful, and are known from Denazinemys nodosa, Basilemys nobilis, Neurankylus baueri, Plastomenus robustus and Thescelus hemispherica. Unidentified taxa are known, including the crurotarsan ?Leidyosuchus, and the theropods ?Struthiomimus, Troodontidae and Tyrannosauridae. The beginning of the Kirtland Formation dates to 74.55 ± 0.22 mya, with the formation ending at around 73.05 ± 0.25 mya.
Utah
Argon-argon radiometric dating indicates that the Kaiparowits Formation was deposited between 76.6 and 74.5 million years ago, during the Campanian age of the Late Cretaceous period. During the Late Cretaceous period, the site of the Kaiparowits Formation was located near the western shore of the Western Interior Seaway, a large inland sea that split North America into two landmasses, Laramidia to the west and Appalachia to the east. The plateau where dinosaurs lived was an ancient floodplain dominated by large channels and abundant wetland peat swamps, ponds and lakes, and was bordered by highlands. The climate was wet and humid, and supported an abundant and diverse range of organisms. This formation contains one of the best and most continuous records of Late Cretaceous terrestrial life in the world.
Parasaurolophus shared its paleoenvironment with other dinosaurs, such as dromaeosaurid theropods, the troodontid Talos sampsoni, ornithomimids like Ornithomimus velox, tyrannosaurids like Teratophoneus, armored ankylosaurids, the duckbilled hadrosaur Gryposaurus monumentensis, the ceratopsians Utahceratops gettyi, Nasutoceratops titusi and Kosmoceratops richardsoni and the oviraptorosaurian Hagryphus giganteus. Paleofauna present in the Kaiparowits Formation included chondrichthyans (sharks and rays), frogs, salamanders, turtles, lizards and crocodilians like the apex predator Deinosuchus. A variety of early mammals were present including multituberculates, marsupials, and insectivorans.
| Biology and health sciences | Ornitischians | Animals |
1974142 | https://en.wikipedia.org/wiki/Depletion%20region | Depletion region | In semiconductor physics, the depletion region, also called depletion layer, depletion zone, junction region, space charge region, or space charge layer, is an insulating region within a conductive, doped semiconductor material where the mobile charge carriers have diffused away, or been forced away by an electric field. The only elements left in the depletion region are ionized donor or acceptor impurities. This region of uncovered positive and negative ions is called the depletion region due to the depletion of carriers in this region, leaving none to carry a current. Understanding the depletion region is key to explaining modern semiconductor electronics: diodes, bipolar junction transistors, field-effect transistors, and variable capacitance diodes all rely on depletion region phenomena.
Formation in a p–n junction
A depletion region forms instantaneously across a p–n junction. It is most easily described when the junction is in thermal equilibrium or in a steady state: in both of these cases the properties of the system do not vary in time; they are in dynamic equilibrium.
Electrons and holes diffuse into regions with lower concentrations of them, much as ink diffuses into water until it is uniformly distributed. By definition, the N-type semiconductor has an excess of free electrons (in the conduction band) compared to the P-type semiconductor, and the P-type has an excess of holes (in the valence band) compared to the N-type. Therefore, when N-doped and P-doped semiconductors are placed together to form a junction, free electrons in the N-side conduction band migrate (diffuse) into the P-side conduction band, and holes in the P-side valence band migrate into the N-side valence band.
Following transfer, the diffused electrons come into contact with holes and are eliminated by recombination in the P-side. Likewise, the diffused holes are recombined with free electrons so eliminated in the N-side. The net result is that the diffused electrons and holes are gone. In a N-side region near to the junction interface, free electrons in the conduction band are gone due to (1) the diffusion of electrons to the P-side and (2) recombination of electrons to holes that are diffused from the P-side. Holes in a P-side region near to the interface are also gone by a similar reason. As a result, majority charge carriers (free electrons for the N-type semiconductor, and holes for the P-type semiconductor) are depleted in the region around the junction interface, so this region is called the depletion region or depletion zone. Due to the majority charge carrier diffusion described above, the depletion region is charged; the N-side of it is positively charged and the P-side of it is negatively charged. This creates an electric field that provides a force opposing the charge diffusion. When the electric field is sufficiently strong to cease further diffusion of holes and electrons, the depletion region reaches the equilibrium. Integrating the electric field across the depletion region determines what is called the built-in voltage (also called the junction voltage or barrier voltage or contact potential).
Physically speaking, charge transfer in semiconductor devices is from (1) the charge carrier drift by the electric field and (2) the charge carrier diffusion due to the spatially varying carrier concentration. In the P-side of the depletion region, where holes drift by the electric field with the electrical conductivity σ and diffuse with the diffusion constant D, the net current density is given by
,
where is the electric field, e is the elementary charge (1.6×10−19 coulomb), and p is the hole density (number per unit volume). The electric field makes holes drift along the field direction, and for diffusion holes move in the direction of decreasing concentration, so for holes a negative current results for a positive density gradient. (If the carriers are electrons, the hole density p is replaced by the electron density n with negative sign; in some cases, both electrons and holes must be included.) When the two current components balance, as in the p–n junction depletion region at dynamic equilibrium, the current is zero due to the Einstein relation, which relates D to σ.
Forward bias
Forward bias (applying a positive voltage to the P-side with respect to the N-side) narrows the depletion region and lowers the barrier to carrier injection (shown in the figure to the right). In more detail, majority carriers get some energy from the bias field, enabling them to go into the region and neutralize opposite charges. The more bias the more neutralization (or screening of ions in the region) occurs. The carriers can be recombined to the ions but thermal energy immediately makes recombined carriers transition back as Fermi energy is in proximity. When bias is strong enough that the depletion region becomes very thin, the diffusion component of the current (through the junction interface) greatly increases and the drift component decreases. In this case, the net current flows from the P-side to the N-side. The carrier density is large (it varies exponentially with the applied bias voltage), making the junction conductive and allowing a large forward current. The mathematical description of the current is provided by the Shockley diode equation. The low current conducted under reverse bias and the large current under forward bias is an example of rectification.
Reverse bias
Under reverse bias (applying a negative voltage to the P-side with respect to the N-side), the potential drop (i.e., voltage) across the depletion region increases. Essentially, majority carriers are pushed away from the junction, leaving behind more charged ions. Thus the depletion region is widened and its field becomes stronger, which increases the drift component of current (through the junction interface) and decreases the diffusion component. In this case, the net current flows from the N-side to the P-side. The carrier density (mostly, minority carriers) is small and only a very small reverse saturation current flows.
Determining the depletion layer width
From a full depletion analysis as shown in figure 2, the charge would be approximated with a sudden drop at its limit points which in reality is gradual and is explained by Poisson's equation. The amount of flux density would then be
where and are the amount of negative and positive charge respectively, and are the distance for negative and positive charge respectively with zero at the center, and are the amount of acceptor and donor atoms respectively and is the electron charge.
Taking the integral of the flux density with respect to distance to determine electric field (i.e. Gauss's law) creates the second graph as shown in figure 2:
where is the permittivity of the substance. Integrating electric field with respect to distance determines the electric potential . This would also equal to the built in voltage as shown in Figure 2.
The final equation would then be arranged so that the function of depletion layer width would be dependent on the electric potential .
In summary, and are the negative and positive depletion layer width respectively with respect to the center, and are the concentration of acceptor and donor atoms respectively, is the electron charge and is the built-in voltage, which is usually the independent variable.
Formation in an MOS capacitor
Another example of a depletion region occurs in the MOS capacitor. It is shown in the figure to the right, for a P-type substrate. Supposing that the semiconductor initially is charge neutral, with the charge due to holes exactly balanced by the negative charge due to acceptor doping impurities. If a positive voltage now is applied to the gate, which is done by introducing positive charge Q to the gate, then some positively charged holes in the semiconductor nearest the gate are repelled by the positive charge on the gate, and exit the device through the bottom contact. They leave behind a depleted region that is insulating because no mobile holes remain; only the immobile, negatively charged acceptor impurities. The greater the positive charge placed on the gate, the more positive the applied gate voltage, and the more holes that leave the semiconductor surface, enlarging the depletion region. (In this device there is a limit to how wide the depletion width may become. It is set by the onset of an inversion layer of carriers in a thin layer, or channel, near the surface. The above discussion applies for positive voltages low enough that an inversion layer does not form.)
If the gate material is polysilicon of opposite type to the bulk semiconductor, then a spontaneous depletion region forms if the gate is electrically shorted to the substrate, in much the same manner as described for the p–n junction above. For more on this, see polysilicon depletion effect.
The principle of charge neutrality says the sum of positive charges must equal the sum of negative charges:
where n and p are the number of free electrons and holes, and and are the number of ionized donors and acceptors "per unit of length", respectively. In this way, both and can be viewed as doping spatial densities. If we assume full ionization and that , then:
.
where and are depletion widths in the p and n semiconductor, respectively. This condition ensures that the net negative acceptor charge exactly balances the net positive donor charge. The total depletion width in this case is the sum . A full derivation for the depletion width is presented in reference. This derivation is based on solving the Poisson equation in one dimension – the dimension normal to the metallurgical junction. The electric field is zero outside of the depletion width (seen in above figure) and therefore Gauss's law implies that the charge density in each region balance – as shown by the first equation in this sub-section. Treating each region separately and substituting the charge density for each region into the Poisson equation eventually leads to a result for the depletion width. This result for the depletion width is:
where is the relative dielectric permittivity of the semiconductor, is the built-in voltage, and is the applied bias. The depletion region is not symmetrically split between the n and p regions - it will tend towards the lightly doped side. A more complete analysis would take into account that there are still some carriers near the edges of the depletion region. This leads to an additional -2kT/q term in the last set of parentheses above.
Depletion width in MOS capacitor
As in p–n junctions, the governing principle here is charge neutrality. Let us assume a P-type substrate. If positive charge Q is placed on gate with area A, then holes are depleted to a depth w exposing sufficient negative acceptors to exactly balance the gate charge. Supposing the dopant density to be acceptors per unit volume, then charge neutrality requires the depletion width w to satisfy the relationship:
If the depletion width becomes wide enough, then electrons appear in a very thin layer at the semiconductor-oxide interface, called an inversion layer because they are oppositely charged to the holes that prevail in a P-type material. When an inversion layer forms, the depletion width ceases to expand with increase in gate charge Q. In this case, neutrality is achieved by attracting more electrons into the inversion layer. In the MOSFET, this inversion layer is referred to as the channel.
Electric field in depletion layer and band bending
Associated with the depletion layer is an effect known as band bending. This effect occurs because the electric field in the depletion layer varies linearly in space from its (maximum) value at the gate to zero at the edge of the depletion width:
where = 8.854×10−12 F/m, F is the farad and m is the meter. This linearly-varying electric field leads to an electrical potential that varies quadratically in space. The energy levels, or energy bands, bend in response to this potential.
| Physical sciences | Electrical circuits | Physics |
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