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1975073 | https://en.wikipedia.org/wiki/Aquatic%20ecosystem | Aquatic ecosystem | An aquatic ecosystem is an ecosystem found in and around a body of water, in contrast to land-based terrestrial ecosystems. Aquatic ecosystems contain communities of organisms—aquatic life—that are dependent on each other and on their environment. The two main types of aquatic ecosystems are marine ecosystems and freshwater ecosystems. Freshwater ecosystems may be lentic (slow moving water, including pools, ponds, and lakes); lotic (faster moving water, for example streams and rivers); and wetlands (areas where the soil is saturated or inundated for at least part of the time).
Types
Marine ecosystems
Marine coastal ecosystem
Marine surface ecosystem
Freshwater ecosystems
Lentic ecosystem (lakes)
Lotic ecosystem (rivers)
Wetlands
Functions
Aquatic ecosystems perform many important environmental functions. For example, they recycle nutrients, purify water, attenuate floods, recharge ground water and provide habitats for wildlife. The biota of an aquatic ecosystem contribute to its self-purification, most notably microorganisms, phytoplankton, higher plants, invertebrates, fish, bacteria, protists, aquatic fungi, and more. These organisms are actively involved in multiple self-purification processes, including organic matter destruction and water filtration. It is crucial that aquatic ecosystems are reliably self-maintained, as they also provide habitats for species that reside in them.
In addition to environmental functions, aquatic ecosystems are also used for human recreation, and are very important to the tourism industry, especially in coastal regions. They are also used for religious purposes, such as the worshipping of the Jordan River by Christians, and educational purposes, such as the usage of lakes for ecological study.
Biotic characteristics (living components)
The biotic characteristics are mainly determined by the organisms that occur. For example, wetland plants may produce dense canopies that cover large areas of sediment—or snails or geese may graze the vegetation leaving large mud flats. Aquatic environments have relatively low oxygen levels, forcing adaptation by the organisms found there. For example, many wetland plants must produce aerenchyma to carry oxygen to roots. Other biotic characteristics are more subtle and difficult to measure, such as the relative importance of competition, mutualism or predation. There are a growing number of cases where predation by coastal herbivores including snails, geese and mammals appears to be a dominant biotic factor.
Autotrophic organisms
Autotrophic organisms are producers that generate organic compounds from inorganic material. Algae use solar energy to generate biomass from carbon dioxide and are possibly the most important autotrophic organisms in aquatic environments. The more shallow the water, the greater the biomass contribution from rooted and floating vascular plants. These two sources combine to produce the extraordinary production of estuaries and wetlands, as this autotrophic biomass is converted into fish, birds, amphibians and other aquatic species.
Chemosynthetic bacteria are found in benthic marine ecosystems. These organisms are able to feed on hydrogen sulfide in water that comes from volcanic vents. Great concentrations of animals that feed on these bacteria are found around volcanic vents. For example, there are giant tube worms (Riftia pachyptila) 1.5 m in length and clams (Calyptogena magnifica) 30 cm long.
Heterotrophic organisms
Heterotrophic organisms consume autotrophic organisms and use the organic compounds in their bodies as energy sources and as raw materials to create their own biomass.
Euryhaline organisms are salt tolerant and can survive in marine ecosystems, while stenohaline or salt intolerant species can only live in freshwater environments.
Abiotic characteristics (non-living components)
An ecosystem is composed of biotic communities that are structured by biological interactions and abiotic environmental factors. Some of the important abiotic environmental factors of aquatic ecosystems include substrate type, water depth, nutrient levels, temperature, salinity, and flow. It is often difficult to determine the relative importance of these factors without rather large experiments. There may be complicated feedback loops. For example, sediment may determine the presence of aquatic plants, but aquatic plants may also trap sediment, and add to the sediment through peat.
The amount of dissolved oxygen in a water body is frequently the key substance in determining the extent and kinds of organic life in the water body. Fish need dissolved oxygen to survive, although their tolerance to low oxygen varies among species; in extreme cases of low oxygen, some fish even resort to air gulping. Plants often have to produce aerenchyma, while the shape and size of leaves may also be altered. Conversely, oxygen is fatal to many kinds of anaerobic bacteria.
Nutrient levels are important in controlling the abundance of many species of algae. The relative abundance of nitrogen and phosphorus can in effect determine which species of algae come to dominate. Algae are a very important source of food for aquatic life, but at the same time, if they become over-abundant, they can cause declines in fish when they decay. Similar over-abundance of algae in coastal environments such as the Gulf of Mexico produces, upon decay, a hypoxic region of water known as a dead zone.
The salinity of the water body is also a determining factor in the kinds of species found in the water body. Organisms in marine ecosystems tolerate salinity, while many freshwater organisms are intolerant of salt. The degree of salinity in an estuary or delta is an important control upon the type of wetland (fresh, intermediate, or brackish), and the associated animal species. Dams built upstream may reduce spring flooding, and reduce sediment accretion, and may therefore lead to saltwater intrusion in coastal wetlands.
Freshwater used for irrigation purposes often absorbs levels of salt that are harmful to freshwater organisms.
Threats
The health of an aquatic ecosystem is degraded when the ecosystem's ability to absorb a stress has been exceeded. A stress on an aquatic ecosystem can be a result of physical, chemical or biological alterations to the environment. Physical alterations include changes in water temperature, water flow and light availability. Chemical alterations include changes in the loading rates of biostimulatory nutrients, oxygen-consuming materials, and toxins. Biological alterations include over-harvesting of commercial species and the introduction of exotic species. Human populations can impose excessive stresses on aquatic ecosystems. Climate change driven by anthropogenic activities can harm aquatic ecosystems by disrupting current distribution patterns of plants and animals. It has negatively impacted deep sea biodiversity, coastal fish diversity, crustaceans, coral reefs, and other biotic components of these ecosystems. Human-made aquatic ecosystems, such as ditches, aquaculture ponds, and irrigation channels, may also cause harm to naturally occurring ecosystems by trading off biodiversity with their intended purposes. For instance, ditches are primarily used for drainage, but their presence also negatively affects biodiversity.
There are many examples of excessive stresses with negative consequences. The environmental history of the Great Lakes of North America illustrates this problem, particularly how multiple stresses, such as water pollution, over-harvesting and invasive species can combine. The Norfolk Broadlands in England illustrate similar decline with pollution and invasive species. Lake Pontchartrain along the Gulf of Mexico illustrates the negative effects of different stresses including levee construction, logging of swamps, invasive species and salt water intrusion.
| Physical sciences | Water: General | Earth science |
1975220 | https://en.wikipedia.org/wiki/Labia | Labia | The labia are the major externally visible structures of the vulva. In humans and other primates, there are two pairs of labia: the labia majora (outer lips) are large and thick folds of skin that cover the vulva's other parts, while the labia minora (inner lips) are the folds of skin between the outer labia that surround and protect the urethral and vaginal openings, as well as the glans clitoris.
In other mammals, the labia majora are not present and the labia minora are instead referred to as the labia vulvae.
Etymology
Labium (plural labia) is a Latin-derived term meaning "lip". Labium and its derivatives (including labial, labrum) are used to describe any lip-like structure, but in the English language, labia often specifically refers to parts of the vulva.
Structure
The labia majora are lip-like structures consisting mostly of skin and adipose (fatty) tissue, which extend on either side of the vulva to form the pudendal cleft through the middle. They enclose and protect the other tissues of the vulva. The labia majora often have a plump appearance, and are thicker towards the anterior. The anterior junction of the labia majora is called the anterior commissure, which is below the mons pubis and above the clitoris. To the posterior, the labia majora join at the posterior commissure, which is above the perineum and below the frenulum of the labia minora. The grooves between the labia majora and labia minora are known as the interlabial sulci or interlabial folds.
The labia minora are two soft folds of fat-free, hairless skin between the labia majora. They enclose and protect the vulvar vestibule, urethra and vagina. The upper portion of each of the labia minora splits to join both the clitoral glans, and the clitoral hood. The labia minora meet posteriorly at the frenulum of the labia minora (also known as the fourchette), which is a fold of skin below the vaginal orifice. The fourchette is more prominent in younger women, and often recedes after sexual activity and childbirth.
When standing or with the legs together, the labia majora usually entirely or partially cover the moist, sensitive inner surfaces of the vulva, which indirectly protects the vagina and urethra, much like the lips protect the mouth. The outer surface of the labia majora is pigmented skin, and develops pubic hair during puberty. The inner surface of the labia majora is smooth, hairless skin, which resembles a mucous membrane, and is only visible when the labia majora and labia minora are drawn apart.
Both the inner and outer surfaces of the labia majora contain sebaceous glands (oil glands), apocrine sweat glands, and eccrine sweat glands. The labia majora have fewer superficial nerve endings than the rest of the vulva, but the skin is highly vascularized. The internal surface of the labia minora is a thin moist skin, with the appearance of a mucous membrane. They contain many sebaceous glands, and occasionally have eccrine sweat glands. The labia minora have many sensory nerve endings, and have a core of erectile tissue.
Diversity
The color, size, length, and shape of the inner labia can vary extensively from woman to woman. In some women, the labia minora are almost non-existent, and in others, they can be fleshy and protuberant. They can range in color from a light pink to brownish black, and texturally can vary between smooth and very rugose.
Embryonic development and changes over time
The urogenital folds form the labia minora while the labioscrotal swellings become the labia majora.
The genital tissues are greatly influenced by natural fluctuations in hormone levels, which lead to changes in labia size, appearance, and elasticity at various life stages. At birth, the labia minora are well-developed, and the labia majora appear plump due to being exposed to maternal hormones in the womb. The labia majora have the same color as the surrounding skin. Labial adhesions can occur between the ages of 3 months and 2 years, and may make the vulva look flat. These adhesions are not usually a cause for concern, and usually disappear without treatment. Treatment options may include estrogen cream, manual separation with local anesthesia, or surgical separation under sedation.
During early childhood, the labia majora look flat and smooth because of decreasing levels of body fat, and the diminished effects of maternal hormones. The labia minora become less prominent.
During puberty, increased hormone levels often significantly change the appearance of the labia. The labia minora become more elastic, prominent, and wrinkled. The labia majora regain fat, and begin growing pubic hair close to the pudendal cleft. Hair is initially sparse and straight, but gradually becomes darker, denser, and curlier as growth spreads outward and upward toward the thighs and mons pubis. At the end of puberty, pubic hair will be coarse, curly, and fairly thick. The patch of pubic hair covering the genitals will eventually often form a triangle shape.
By adulthood, the outer surface of the labia majora may be darker than the surrounding skin, and may have wrinkles. During the reproductive years, if a woman delivers a child, the fourchette will flatten. Pregnancy may cause the labia minora to darken in color.
Later in life, the labia majora once again gradually lose fat, becoming flatter and more wrinkled, and pubic hair turns grey. Following menopause, falling hormone levels cause further changes to the labia. The labia minora atrophy, making them become less elastic, and pubic hair on the labia majora becomes more sparse.
Sexual arousal and response
The labia are an erogenous zone. In particular, the labia minora are sexually responsive, and sensitivity varies greatly. In some, they are so sensitive that anything other than light touch may be uncomfortable, whereas stimulation may elicit no sexual response in others. The labia may be sexually stimulated as part of masturbation or with a sex partner, such as by fingering or oral sex. Moving the labia minora can also stimulate the extremely sensitive clitoris.
During sexual arousal, the labia majora swell due to increased blood flow to the region, and slightly draw back, revealing the inner labia. The labia minora become engorged with blood, causing them to expand in diameter by two to three times, and darken or redden in color. Because pregnancy and childbirth increase genital vascularity, the inner and outer labia will engorge faster in women who have had children.
After a period of sexual stimulation, the labia minora will become further engorged with blood approximately 30 seconds to 3 minutes before orgasm, causing them to redden further. In those who have had children, the labia majora may also swell significantly during this period, becoming dark red. Continued stimulation can result in an orgasm, and the orgasmic contractions help remove blood trapped in the inner and outer labia, as well as the clitoris and other parts of the vulva, which causes pleasurable orgasmic sensations.
Following orgasm or when no longer sexually aroused, the labia gradually return to their unaroused state. The labia minora return to their original color within 2 minutes, and engorgement dissipates in about 5 to 10 minutes. The labia majora return to their pre-arousal state in approximately 1 hour.
Society and culture
In many cultures and locations all over the world, the labia, as part of the genitalia, are considered private, or intimate parts, whose exposure (especially in public) is governed by fairly strict socio-cultural mores. In many cases, public exposure is limited, and often prohibited by law.
Views on pubic hair differ between people and between cultures. Some women prefer the look or feel of pubic hair, while others may choose to remove some or all of it. Temporary methods of removal include shaving, trimming, waxing, sugaring and depilatory products while permanent hair removal can be accomplished using electrolysis or laser hair removal. In Korea, pubic hair is considered a sign of fertility, leading some women to have pubic hair transplants.
Some women are self-conscious about the size, color or asymmetry of their labia. Viewing pornography may influence a woman's view of her genitals. Models in pornography frequently have small or non-existent labia minora, and images are often airbrushed, so pornographic images do not depict the full range of natural variations of the vulva. This can lead viewers of pornography to have unrealistic expectations about how the labia should look. Similar to how some women develop self-esteem issues from comparing their faces and bodies to airbrushed models in magazines, women who compare their vulvas to idealized pornographic images may believe their own labia are abnormal. This can have a negative impact on a woman's life, since genital self-consciousness makes it more difficult to enjoy sexual activity, see a gynecologist, or perform a genital self-examination. Developing an awareness for how much the labia truly differ between individuals may help to overcome this self-consciousness.
In several countries in Africa and Asia, the vulva is routinely altered or removed for reasons related to ideas about tradition, purity, hygiene and aesthetics. Known as female genital mutilation, the procedures include clitoridectomy and so-called "pharaonic circumcision," whereby the inner and outer labia are removed and the vulva is sewn shut. FGM is mostly outlawed around the world, even in countries where the practice is widespread.
Labiaplasty is a controversial plastic surgery procedure that involves the creation or reshaping of the labia. Labia piercing is a cosmetic piercing, usually with a special needle under sterile conditions, of the inner or outer labia. Jewelry is worn in the resulting opening.
Additional images
| Biology and health sciences | Human anatomy | Health |
1975821 | https://en.wikipedia.org/wiki/Skew%20lines | Skew lines | In three-dimensional geometry, skew lines are two lines that do not intersect and are not parallel. A simple example of a pair of skew lines is the pair of lines through opposite edges of a regular tetrahedron. Two lines that both lie in the same plane must either cross each other or be parallel, so skew lines can exist only in three or more dimensions. Two lines are skew if and only if they are not coplanar.
General position
If four points are chosen at random uniformly within a unit cube, they will almost surely define a pair of skew lines. After the first three points have been chosen, the fourth point will define a non-skew line if, and only if, it is coplanar with the first three points. However, the plane through the first three points forms a subset of measure zero of the cube, and the probability that the fourth point lies on this plane is zero. If it does not, the lines defined by the points will be skew.
Similarly, in three-dimensional space a very small perturbation of any two parallel or intersecting lines will almost certainly turn them into skew lines. Therefore, any four points in general position always form skew lines.
In this sense, skew lines are the "usual" case, and parallel or intersecting lines are special cases.
Formulas
Testing for skewness
Nearest points
Expressing the two lines as vectors:
The cross product of and is perpendicular to the lines.
The plane formed by the translations of Line 2 along contains the point and is perpendicular to .
Therefore, the intersecting point of Line 1 with the above-mentioned plane, which is also the point on Line 1 that is nearest to Line 2 is given by
Similarly, the point on Line 2 nearest to Line 1 is given by (where )
Distance
The nearest points and form the shortest line segment joining Line 1 and Line 2:
The distance between nearest points in two skew lines may also be expressed using other vectors:
Here the 1×3 vector represents an arbitrary point on the line through particular point with representing the direction of the line and with the value of the real number determining where the point is on the line, and similarly for arbitrary point on the line through particular point in direction .
The cross product of b and d is perpendicular to the lines, as is the unit vector
The perpendicular distance between the lines is then
(if |b × d| is zero the lines are parallel and this method cannot be used).
More than two lines
Configurations
A configuration of skew lines is a set of lines in which all pairs are skew. Two configurations are said to be isotopic if it is possible to continuously transform one configuration into the other, maintaining throughout the transformation the invariant that all pairs of lines remain skew. Any two configurations of two lines are easily seen to be isotopic, and configurations of the same number of lines in dimensions higher than three are always isotopic, but there exist multiple non-isotopic configurations of three or more lines in three dimensions. The number of nonisotopic configurations of n lines in R3, starting at n = 1, is
1, 1, 2, 3, 7, 19, 74, ... .
Ruled surfaces
An affine transformation of this ruled surface produces a surface which in general has an elliptical cross-section rather than the circular cross-section produced by rotating L around L'; such surfaces are also called hyperboloids of one sheet, and again are ruled by two families of mutually skew lines. A third type of ruled surface is the hyperbolic paraboloid. Like the hyperboloid of one sheet, the hyperbolic paraboloid has two families of skew lines; in each of the two families the lines are parallel to a common plane although not to each other. Any three skew lines in R3 lie on exactly one ruled surface of one of these types.
Gallucci's theorem
If three skew lines all meet three other skew lines, any transversal of the first set of three meets any transversal of the second set.
Skew flats in higher dimensions
In higher-dimensional space, a flat of dimension k is referred to as a k-flat. Thus, a line may also be called a 1-flat.
Generalizing the concept of skew lines to d-dimensional space, an i-flat and a j-flat may be skew if
. As with lines in 3-space, skew flats are those that are neither parallel nor intersect.
In affine d-space, two flats of any dimension may be parallel.
However, in projective space, parallelism does not exist; two flats must either intersect or be skew.
Let be the set of points on an i-flat, and let be the set of points on a j-flat.
In projective d-space, if then the intersection of and must contain a (i+j−d)-flat. (A 0-flat is a point.)
In either geometry, if and intersect at a k-flat, for , then the points of determine a (i+j−k)-flat.
| Mathematics | Three-dimensional space | null |
1975956 | https://en.wikipedia.org/wiki/Electric%20potential%20energy | Electric potential energy | Electric potential energy is a potential energy (measured in joules) that results from conservative Coulomb forces and is associated with the configuration of a particular set of point charges within a defined system. An object may be said to have electric potential energy by virtue of either its own electric charge or its relative position to other electrically charged objects.
The term "electric potential energy" is used to describe the potential energy in systems with time-variant electric fields, while the term "electrostatic potential energy" is used to describe the potential energy in systems with time-invariant electric fields.
Definition
The electric potential energy of a system of point charges is defined as the work required to assemble this system of charges by bringing them close together, as in the system from an infinite distance. Alternatively, the electric potential energy of any given charge or system of charges is termed as the total work done by an external agent in bringing the charge or the system of charges from infinity to the present configuration without undergoing any acceleration.
The electrostatic potential energy can also be defined from the electric potential as follows:
Units
The SI unit of electric potential energy is joule (named after the English physicist James Prescott Joule). In the CGS system the erg is the unit of energy, being equal to 10−7 Joules. Also electronvolts may be used, 1 eV = 1.602×10−19 Joules.
Electrostatic potential energy of one point charge
One point charge q in the presence of another point charge Q
The electrostatic potential energy, UE, of one point charge q at position r in the presence of a point charge Q, taking an infinite separation between the charges as the reference position, is:
where r is the distance between the point charges q and Q, and q and Q are the charges (not the absolute values of the charges—i.e., an electron would have a negative value of charge when placed in the formula). The following outline of proof states the derivation from the definition of electric potential energy and Coulomb's law to this formula.
One point charge q in the presence of n point charges Qi
The electrostatic potential energy, UE, of one point charge q in the presence of n point charges Qi, taking an infinite separation between the charges as the reference position, is:
where ri is the distance between the point charges q and Qi, and q and Qi are the assigned values of the charges.
Electrostatic potential energy stored in a system of point charges
The electrostatic potential energy UE stored in a system of N charges q1, q2, …, qN at positions r1, r2, …, rN respectively, is:
where, for each i value, V(ri) is the electrostatic potential due to all point charges except the one at ri, and is equal to:
where rij is the distance between qi and qj.
Energy stored in a system of one point charge
The electrostatic potential energy of a system containing only one point charge is zero, as there are no other sources of electrostatic force against which an external agent must do work in moving the point charge from infinity to its final location.
A common question arises concerning the interaction of a point charge with its own electrostatic potential. Since this interaction doesn't act to move the point charge itself, it doesn't contribute to the stored energy of the system.
Energy stored in a system of two point charges
Consider bringing a point charge, q, into its final position near a point charge, Q1. The electric potential V(r) due to Q1 is
Hence we obtain, the electrostatic potential energy of q in the potential of Q1 as
where r1 is the separation between the two point charges.
Energy stored in a system of three point charges
The electrostatic potential energy of a system of three charges should not be confused with the electrostatic potential energy of Q1 due to two charges Q2 and Q3, because the latter doesn't include the electrostatic potential energy of the system of the two charges Q2 and Q3.
The electrostatic potential energy stored in the system of three charges is:
Energy stored in an electrostatic field distribution in vacuum
The energy density, or energy per unit volume, , of the electrostatic field of a continuous charge distribution is:
Energy stored in electronic elements
Some elements in a circuit can convert energy from one form to another. For example, a resistor converts electrical energy to heat. This is known as the Joule effect. A capacitor stores it in its electric field. The total electrostatic potential energy stored in a capacitor is given by
where C is the capacitance, V is the electric potential difference, and Q the charge stored in the capacitor.
The total electrostatic potential energy may also be expressed in terms of the electric field in the form
where is the electric displacement field within a dielectric material and integration is over the entire volume of the dielectric.
The total electrostatic potential energy stored within a charged dielectric may also be expressed in terms of a continuous volume charge, ,
where integration is over the entire volume of the dielectric.
These latter two expressions are valid only for cases when the smallest increment of charge is zero () such as dielectrics in the presence of metallic electrodes or dielectrics containing many charges.
Note that a virtual experiment based on the energy transfert between capacitor plates reveals that an additional term should be taken into account when dealing with semiconductors for instance . While this extra energy cancels when dealing with insulators, the derivation predicts that it cannot be ignored as it may exceed the polarization energy.
| Physical sciences | Electrostatics | Physics |
1976413 | https://en.wikipedia.org/wiki/Divergent%20evolution | Divergent evolution | Divergent evolution or divergent selection is the accumulation of differences between closely related populations within a species, sometimes leading to speciation. Divergent evolution is typically exhibited when two populations become separated by a geographic barrier (such as in allopatric or peripatric speciation) and experience different selective pressures that cause adaptations. After many generations and continual evolution, the populations become less able to interbreed with one another. The American naturalist J. T. Gulick (1832–1923) was the first to use the term "divergent evolution", with its use becoming widespread in modern evolutionary literature. Examples of divergence in nature are the adaptive radiation of the finches of the Galápagos, changes in mobbing behavior of the kittiwake, and the evolution of the modern-day dog from the wolf.
The term can also be applied in molecular evolution, such as to proteins that derive from homologous genes. Both orthologous genes (resulting from a speciation event) and paralogous genes (resulting from gene duplication) can illustrate divergent evolution. Through gene duplication, it is possible for divergent evolution to occur between two genes within a species. Similarities between species that have diverged are due to their common origin, so such similarities are homologies.
Causes
Animals undergo divergent evolution for a number of reasons linked to changes in environmental or social pressures. This could include changes in the environment, such access to food and shelter. It could also result from changes in predators, such as new adaptations, an increase or decrease in number of active predators, or the introduction of new predators. Divergent evolution can also be a result of mating pressures such as increased competition for mates or selective breeding by humans.
Distinctions
Divergent evolution is a type of evolution and is distinct from convergent evolution and parallel evolution, although it does share similarities with the other types of evolution.
Divergent versus convergent evolution
Convergent evolution is the development of analogous structures that occurs in different species as a result of those two species facing similar environmental pressures and adapting in similar ways. It differs from divergent evolution as the species involved do not descend from a closely related common ancestor and the traits accumulated are similar. An example of convergent evolution is the development of flight in birds, bats, and insects, all of which are not closely related but share analogous structures allowing for flight.
Divergent versus parallel evolution
Parallel evolution is the development of a similar trait in species descending from a common ancestor. It is comparable to divergent evolution in that the species are descend from a common ancestor, but the traits accumulated are similar due to similar environmental pressures while in divergent evolution the traits accumulated are different. An example of parallel evolution is that certain arboreal frog species, 'flying' frogs, in both Old World families and New World families, have developed the ability of gliding flight. They have "enlarged hands and feet, full webbing between all fingers and toes, lateral skin flaps on the arms and legs, and reduced weight per snout-vent length".
Darwin's finches
One of the first recorded examples of divergent evolution is the case of Darwin's Finches. During Darwin's travels to the Galápagos Islands, he discovered several different species of finch, living on the different islands. Darwin observed that the finches had different beaks specialized for that species of finches' diet. Some finches had short beaks for eating nuts and seeds, other finches had long thin beaks for eating insects, and others had beaks specialized for eating cacti and other plants. He concluded that the finches evolved from a shared common ancestor that lived on the islands, and due to geographic isolation, evolved to fill the particular niche on each of the islands. This is supported by modern day genomic sequencing.
Divergent evolution in dogs
Another example of divergent evolution is the origin of the domestic dog and the modern wolf, who both shared a common ancestor. Comparing the anatomy of dogs and wolves supports this claim as they have similar body shape, skull size, and limb formation. This is even more obvious in some species of dogs, such as malamutes and huskies, who appear even more physically and behaviorally similar. There is a divergent genomic sequence of the mitochondrial DNA of wolves and dogs dated to over 100,000 years ago, which further supports the theory that dogs and wolves have diverged from shared ancestry.
Divergent evolution in kittiwakes
Another example of divergent evolution is the behavioral changes in the kittiwake as opposed to other species of gulls. Ancestorial and other modern-day species of gulls exhibit a mobbing behavior in order to protect their young due the nesting at ground-level where they are susceptible to predators. As a result of migration and environmental changes, the kittiwake nest solely on cliff faces. As a result, their young are protected from predatory reptiles, mammals, and birds who struggle with the climb and cliff-face weather conditions, and they do not exhibit this mobbing behavior.
Divergent evolution in cacti
Another example of divergent evolution is the split forming the Cactaceae family approximately dated in the late Miocene. Due to increase in arid climates, following the Eocene–Oligocene event, these ancestral plants evolved to survive in the new climates. Cacti evolved to have areoles, succulent stems, and some have light leaves, with the ability to store water for up to months. The plants they diverged from either went extinct leaving little in the fossil record or migrated surviving in less arid climates.
| Biology and health sciences | Basics_4 | Biology |
1976615 | https://en.wikipedia.org/wiki/Metanephrops%20challengeri | Metanephrops challengeri | Metanephrops challengeri (commonly known as the New Zealand lobster or New Zealand scampi) is a species of slim, pink lobster that lives around the coast of New Zealand. It is typically long and weighs around . The carapace and abdomen are smooth, and adults are white with pink and brown markings and a conspicuous pair of long, slim claws. M. challengeri lives in burrows at depths of in a variety of sediments. Although individuals can live for up to 15 years, the species shows low fecundity, where small numbers of larvae hatch at an advanced stage.
M. challengeri is a significant prey item for ling, as well as being an important fishery species for human consumption; trawlers catch around per year under the limitations of New Zealand's Quota Management System. The species was first collected by the Challenger expedition of 1872–1876, but only described as separate from related species by Heinrich Balss in 1914. Although originally classified in the genus Nephrops, it was moved in 1972 to a new genus, Metanephrops, along with most other species then classified in Nephrops.
Description
Metanephrops challengeri is a slender lobster, typically long, but exceptionally up to , and weighing up to each. Its chelipeds (legs bearing the main chelae, or claws) are long, narrow, and slightly unequal. The second and third pairs of pereiopods also end in small claws, but the fourth and fifth pairs do not. The carapace is smooth, and extends forwards into a long, narrow rostrum, only slightly shorter than the carapace.
Adults are mostly white, but the front half of the rostrum, and the sides of the abdomen, are pink. Bright red bands extend across the base of the rostrum, the posterior edge of the carapace, the chelipeds, and each of the abdominal segments. The dorsal parts of the abdomen are brown, and there are two brown saddles on the dorsal carapace.
M. challengeri is considered to have the most primitive morphology of any species of Metanephrops, having even fewer novelties than the oldest known fossil species, M. rossensis. Its rostrum is longer than that of other species in the thomsoni species group, and the ridge along the midline of the carapace only has two small spines. Unlike some other species of Metanephrops, the carapace is smooth, as are the abdominal tergae, and the chelipeds are covered in fine granules.
Life cycle
Metanephrops challengeri reaches sexual maturity at the age of 3–4 years, and may live up to 15 years in total. Females produce very large eggs in small numbers; they are typically around in diameter, and are blue in colour. The larvae hatch at the zoea stage (equivalent to the third zoea of the Northern Hemisphere species Nephrops norvegicus). The zoea larvae are long, and possess all the appendages of the cephalothorax, including the pereiopods, which are used for swimming, but no pleopods (appendages of the abdomen). This larval stage lasts less than four days, before the young moult into the post-larval stage. The post-larva swims using its pleopods. The post-larva later moults into the adult form. Larvae are rarely seen in the wild, confirming that the development to the bottom-dwelling post-larva is rapid.
Distribution and ecology
Metanephrops challengeri lives around the coasts of New Zealand, including the Chatham Islands, at depths of . It lives in burrows in a variety of "suitable cohesive" sediments, and is a significant prey item for ling (Genypterus blacodes). Lobsters have few parasites, the most important for M. challengeri being the microsporidian Myospora metanephrops. This can cause "destruction of the skeletal and heart muscles of infected lobsters", but its significance for the animals and for the fishing industry remains unclear. When it was described in 2010, M. metanephrops was the first microsporidian to be isolated from a true lobster.
Fisheries
Metanephrops challengeri has been harvested commercially since the 1980s. Between the season of 1988/89 and 1990/91, the amount of scampi caught around New Zealand increased from only to around . Catch limits were introduced in 1990/91, and now is caught annually by trawlers. The fishery is centred on four areas of continental shelf of the submerged continent Zealandia: the Campbell Plateau around the Auckland Islands, Chatham Rise, along the Wairarapa coast, and in the Bay of Plenty.
Most of the fishing vessels used to capture M. challengeri are long, with "double or triple trawl rigs of low headline height". There is considerable variation in the catch per unit effort between different depths, between different geographical areas and between different years. M. challengeri is considered a luxury foodstuff. Most of the catch is exported and as a result, it is rarely seen in restaurants in New Zealand.
Metanephrops challengeri was the subject of a 2003 select committee inquiry in the New Zealand parliament, after allegations of corruption arose against officers of the Ministry of Fisheries. Although the allegations were quashed, the inquiry ruled that preferential treatment had been given to the large fishing company Simunovich Fisheries. In response, the government introduced M. challengeri into their Quota Management System and paid compensation to some fishermen who had a justified grievance. Under QMS, an overall limit of was put in place for M. challengeri in 2011.
Conservation
Metanephrops challengeri is currently listed as Least Concern on the IUCN Red List, due in part to the Quota Management System put in place by the New Zealand government. The species does appear to be declining, however, based both on burrow counts and analyses of catch per unit effort. Estimates of the total population size of M. challengeri vary depending on the methods used. Based on indirect measures, such as burrow counts, there may be as many as 28 million individuals, and the annual catch might represent only 2%–4% of the total population. Using more reliable figures based on those animals seen during surveys, there may be only 2–11 million individuals available to trawlers, and the annual catch may remove 12%–28% of that population. Bycatch from the New Zealand scampi fishery has included the New Zealand sea lion, Phocarctos hookeri, which is considered a vulnerable species by the IUCN.
Taxonomy
Metanephrops challengeri was first described by Heinrich Balss in 1914, under the name Nephrops challengeri. Two specimens had been collected on the Challenger expedition from benthic Globigerina ooze at a depth of , on the Challenger Plateau in the Tasman Sea (). They had been included by Charles Spence Bate in his report on the crustaceans collected by the Challenger expedition, but were not separated from "Nephrops thomsoni" (now Metanephrops thomsoni), which was described by Spence Bate as a new species. Balss recognised that Spence Bate's N. thomsoni covered two species and, restricting the name M. thomsoni to the species containing the type specimens designated by Spence Bate (from the Philippines), created a new species for the species from New Zealand. Balss chose the two specimens seen by Spence Bate to be the type specimens of his new species, Nephrops challengeri. Both were females, and they have been deposited at the Natural History Museum in London.
The species was transferred to a new genus, Metanephrops (along with every other extant species then in Nephrops, except its type species, Nephrops norvegicus) by Richard Jenkins of the University of Adelaide in 1972. Jenkins placed M. challengeri among the "thomsoni group" within the genus Metanephrops, alongside M. thomsoni, M. sibogae, M. boschmai and M. sinensis. Jenkins inferred that this group of species had originated off northern Australia or in Indonesia, and that M. challengeri had reached New Zealand in the late Tertiary and displaced M. motunauensis, which formerly lived there. More recently, findings from molecular phylogenetics suggest that M. challengeri has a basal position in the genus, possibly linked to M. neptunus, and that the genus may have originated at high latitudes in the South Atlantic.
| Biology and health sciences | Crayfishes and lobsters | Animals |
1977416 | https://en.wikipedia.org/wiki/Nile%20monitor | Nile monitor | The Nile monitor (Varanus niloticus) is a large member of the monitor family (Varanidae) found throughout most of Sub-Saharan Africa, particularly in drier regions, and along the Nile River and its tributaries in East Africa. Additionally, there are modern, invasive populations in North America. The population found in West African forests and savannahs is sometimes recognized as a separate species, the West African Nile monitor (V. stellatus). While it is dwarfed by its larger relatives, such as the Komodo dragon, the Asian water monitor or the crocodile monitor, it is still one of the largest lizards in the world, reaching (and even surpassing) Australia’s perentie in size. Other common names include the African small-grain lizard, as well as iguana and various forms derived from it, such as guana, water leguaan or river leguaan (leguan, leguaan, and likkewaan mean monitor lizard in South African English, and can be used interchangeably).
A feral population of Nile monitors (descended from escaped or intentionally-released pets) has become established in several locations in South Florida. In addition to any illegally-released animals, it is speculated that during particularly intense hurricane seasons in Florida, many reptiles potentially escape when their enclosures are damaged or inadvertently unlocked; as Florida has a semi-tropical to tropical climate, many reptiles are housed outdoors, and poorly-secured enclosures may become damaged during bad storms. Along with Nile monitors, Florida is infamous for its feral populations of agamas, Argentine black and white tegus, Burmese pythons, green iguanas, Madagascar giant day geckos, and panther and veiled chameleons, among others. Many of these species are thought to be descendants of hurricane escapees.
Taxonomy
Members of the Nile monitor species group were already well known to Africans in ancient times. For example, they were commonly caught, likely as food, in the Djenné-Djenno culture at least a millennium ago.
The Nile monitor twice was given a scientific name by Carl Linnaeus: First as Lacerta monitor in 1758 in the 10th edition of Systema Naturae, the starting point of zoological nomenclature. He described it again in 1766 as Lacerta nilotica. Despite being older, the name proposed in 1758 is invalid because it was rejected in ICZN opinion 540, making the name of 1766 valid. The genus Varanus was coined in 1820 by Blasius Merrem. Six years later Leopold Fitzinger moved the Nile monitor into this genus as Varanus niloticus, the currently accepted scientific name for the species.
Species complex
As traditionally defined, the Nile monitor is a species complex.
The ornate monitor (V. ornatus) and West African Nile monitor (V. stellatus) were described as species in 1802 and 1803 by François Marie Daudin. In 1942, Robert Mertens moved them both into the Nile monitor (V. niloticus); as synonyms or as a valid subspecies. This was the standard treatment until 1997, when a taxonomic review based on color and morphology indicated that the ornate monitor is distinctive and revalidated it as a separate species from rainforests of West and Central Africa. In 2016, a review based primarily on genetics came to another result. They found that monitors from West African forests and adjacent savannah are distinctive and worthy of recognition as a separate species: the West African Nile monitor (V. stellatus). It is estimated to have split from the others in the Nile monitor complex about 7.7 million years ago, making it older than the split between humans and chimpanzees. In contrast, those in the Central African rainforests are genetically similar to the Nile monitor. This essentially splits the ornate monitor—as defined in 1997—into two: the western being the West African Nile monitor and the eastern (of Central African rainforests) being moved back into the Nile monitor. As the type locality for the ornate monitor is in the Central African country of Cameroon, the scientific name V. ornatus becomes a synonym of V. niloticus. Individuals with the "ornate color pattern" and individuals with the "Nile color pattern" occur in both the West African Nile monitor and the Nile monitor, with the "ornate" appearing to be more frequent in densely forested habitats.
With the West African Nile monitor as a separate species, there are two main clades in the Nile monitor: A widespread clade found throughout much of Southern, Central and East Africa, as well as more locally in coastal West Africa. The other clade includes the monitors of the Sahel (Mali to Ethiopia) and Nile regions. Despite the differences, the Reptile Database maintains both the ornate monitor and West African Nile monitor as synonyms of the Nile monitor, but do note that this broad species definition includes distinctive subpopulations.
Description
The Nile monitor is Africa's longest lizard. They grow from about in length, with the largest specimens attaining . In an average-sized specimen, the snout-to-vent length will be around . In body mass, adults have been reported to vary widely, one study claiming only , others state weights ranging from in big monitors. Variations may be due to age or environmental conditions. Exceptionally large specimens may scale as much as , but this species weighs somewhat less on average than the bulkier rock monitor.
They have muscular bodies, strong legs, and powerful jaws. Their teeth are sharp and pointed in juvenile animals and become blunt and peg-like in adults. They also possess sharp claws used for climbing, digging, defense, or tearing at their prey. Like all monitors, they have forked tongues, with highly developed olfactory properties. The Nile monitor has quite striking, but variable, skin patterns, as they are greyish-brown above with greenish-yellow barring on the tail and large, greenish-yellow rosette-like spots on their backs with a blackish tiny spot in the middle. Their throats and undersides are an ochre-yellow to a creamy-yellow, often with faint barring.
Their nostrils are placed high on their snouts, indicating these animals are very well adapted for an aquatic lifestyle. They are also excellent climbers and quick runners on land. Nile monitors feed on a wide variety of prey items, including fish, frogs and toads (even poisonous ones of the genera Breviceps and Sclerophrys), small reptiles (such as turtles, snakes, lizards, and young crocodiles), birds, rodents, other small mammals (up to domestic cats and young antelopes [Raphicerus]), eggs (including those of crocodiles, agamids, other monitor lizards, and birds), invertebrates (such as beetles, termites, orthopterans, crabs, caterpillars, spiders, millipedes, earthworms, snails, and slugs), carrion, human wastes, and feces.
Distribution and habitat
Nile monitors are native to Sub-Saharan Africa and along the Nile. They are not found in any of the desert regions of Africa (notably Sahara, Kalahari and much of the Horn of Africa), however, they thrive around rivers. Nile monitors were reported to live in and around the Jordan River, Dead Sea, and wadis of the Judaean Desert in Israel until the late 19th century, though they are now extinct in the region.
Invasive species
In Florida in the United States, established breeding populations of Nile monitors have been known to exist in different parts of the state since at least 1990. Genetic studies have shown that these introduced animals are part of the subpopulation that originates from West Africa, and now often is recognized as its own species, the West African Nile monitor. The vast majority of the established breeding population is in Lee County, particularly in the Cape Coral and surrounding regions, including the nearby barrier islands (Sanibel, Captiva, and North Captiva), Pine Island, Fort Myers, and Punta Rassa. Established populations also exist in adjacent Charlotte County, especially on Gasparilla Island. Other areas in Florida with a sizeable number of Nile monitor sightings include Palm Beach County just southwest of West Palm Beach along State Road 80. In July 2008, a Nile monitor was spotted in Homestead, a small city southwest of Miami. Other sightings have been reported near Hollywood, Naranja, and as far south as Key Largo in the Florida Keys. The potential for the established population of Nile monitors in Lee, Charlotte, and other counties in Florida, to negatively impact indigenous crocodilians, such as American alligators (Alligator mississippiensis), and American crocodiles (Crocodylus acutus), is enormous, given that they normally raid crocodile nests, eat eggs, and prey on small crocodiles in Africa. Anecdotal evidence indicates a high rate of disappearance of domestic pets and feral cats in Cape Coral.
In captivity
Nile monitors are often found in the pet trade despite their highly aggressive demeanor and resistance to taming. Juvenile monitors will tail whip as a defensive measure, and as adults, they are capable of inflicting moderate to serious wounds from biting and scratching. Nile monitors require a large cage as juveniles quickly grow when fed a varied diet, and large adults often require custom-built quarters.
There are few lizards less suited to life in captivity than the Nile monitor. Buffrenil (1992) considered that, when fighting for its life, a Nile monitor was a more dangerous adversary than a crocodile of a similar size. Their care presents particular problems on account of the lizards' enormous size and lively dispositions. Very few of the people who buy brightly-coloured baby Nile monitors can be aware that, within a couple of years, their purchase will have turned into an enormous, ferocious carnivore, quite capable of breaking the family cat's neck with a single snap and swallowing it whole.
| Biology and health sciences | Lizards and other Squamata | Animals |
875134 | https://en.wikipedia.org/wiki/Galactic%20halo | Galactic halo | A galactic halo is an extended, roughly spherical component of a galaxy which extends beyond the main, visible component. Several distinct components of a galaxy comprise its halo:
the stellar halo
the galactic corona (hot gas, i.e. a plasma)
the dark matter halo
The distinction between the halo and the main body of the galaxy is clearest in spiral galaxies, where the spherical shape of the halo contrasts with the flat disc. In an elliptical galaxy, there is no sharp transition between the other components of the galaxy and the halo.
A halo can be studied by observing its effect on the passage of light from distant bright objects like quasars that are in line of sight beyond the galaxy in question.
Components of the galactic halo
Stellar halo
The stellar halo is a nearly spherical population of field stars and globular clusters. It surrounds most disk galaxies as well as some elliptical galaxies of type cD. A low amount (about one percent) of a galaxy's stellar mass resides in the stellar halo, meaning its luminosity is much lower than other components of the galaxy.
The Milky Way's stellar halo contains globular clusters, RR Lyrae stars with low metallicity, and subdwarfs. In our stellar halo, stars tend to be old (most are greater than 12 billion years old) and metal-poor, but there are also halo star clusters with observed metal content similar to disk stars. The halo stars of the Milky Way have an observed radial velocity dispersion of about 200 kilometres per second and a low average velocity of rotation of about . Star formation in the stellar halo of the Milky Way ceased long ago.
Galactic corona
A galactic corona is a distribution of gas extending far away from the center of the galaxy. It can be detected by the distinct emission spectrum it gives off, showing the presence of atomic neutral hydrogen (the H I region, pronounced "H-one") and other features detectable by X-ray spectroscopy.
Dark matter halo
The dark matter halo is a theorized distribution of dark matter which extends throughout the galaxy extending far beyond its visible components. The mass of the dark matter halo is far greater than the mass of the other components of the galaxy. Its existence is hypothesized in order to account for the gravitational potential that determines the dynamics of bodies within galaxies. The nature of dark matter halos is an important area in current research in cosmology, in particular its relation to galactic formation and evolution.
The Navarro–Frenk–White profile is a widely accepted density profile of the dark matter halo determined through numerical simulations. It represents the mass density of the dark matter halo as a function of , the distance from the galactic center:
where is a characteristic radius for the model, is the critical density (with being the Hubble constant), and is a dimensionless constant. The invisible halo component cannot extend with this density profile indefinitely, however; this would lead to a diverging integral when calculating mass. It does, however, provide a finite gravitational potential for all . Most measurements that can be made are relatively insensitive to the outer halo's mass distribution. This is a consequence of Newton's laws, which state that if the shape of the halo is spheroidal or elliptical there will be no net gravitational effect from halo mass a distance from the galactic center on an object that is closer to the galactic center than . The only dynamical variable related to the extent of the halo that can be constrained is the escape velocity: the fastest-moving stellar objects still gravitationally bound to the Galaxy can give a lower bound on the mass profile of the outer edges of the dark halo.
Formation of galactic halos
The formation of stellar halos occurs naturally in a cold dark matter model of the universe in which the evolution of systems such as halos occurs from the bottom-up, meaning the large scale structure of galaxies is formed starting with small objects. Halos, which are composed of both baryonic and dark matter, form by merging with each other. Evidence suggests that the formation of galactic halos may also be due to the effects of increased gravity and the presence of primordial black holes. The gas from halo mergers goes toward the formation of the central galactic components, while stars and dark matter remain in the galactic halo.
On the other hand, the halo of the Milky Way Galaxy is thought to derive from the Gaia Sausage.
| Physical sciences | Basics_2 | Astronomy |
875148 | https://en.wikipedia.org/wiki/Evolution%20of%20cetaceans | Evolution of cetaceans | The evolution of cetaceans is thought to have begun in the Indian subcontinent from even-toed ungulates (Artiodactyla) 50 million years ago (mya) and to have proceeded over a period of at least 15 million years. Cetaceans are fully aquatic mammals belonging to the order Artiodactyla and branched off from other artiodactyls around 50 mya. Cetaceans are thought to have evolved during the Eocene (56-34 mya), the second epoch of the present-extending Cenozoic Era. Molecular and morphological analyses suggest Cetacea share a relatively recent closest common ancestor with hippopotami and that they are sister groups. Being mammals, they surface to breathe air; they have five finger bones (even-toed) in their fins; they nurse their young; and, despite their fully aquatic life style, they retain many skeletal features from their terrestrial ancestors. Research conducted in the late 1970s in Pakistan revealed several stages in the transition of cetaceans from land to sea.
The two modern parvorders of cetaceans – Mysticeti (baleen whales) and Odontoceti (toothed whales) – are thought to have separated from each other around 28–33 mya in a second cetacean radiation, the first occurring with the archaeocetes. The adaptation of animal echolocation in toothed whales distinguishes them from fully aquatic archaeocetes and early baleen whales. The presence of baleen in baleen whales occurred gradually, with earlier varieties having very little baleen, and their size is linked to baleen dependence (and subsequent increase in filter feeding).
Early evolution
The aquatic lifestyle of cetaceans first began in the Indian subcontinent from even-toed ungulates 50 million years ago, with this initial stage lasting approximately 4-15 million years. Archaeoceti is an extinct parvorder of Cetacea containing ancient whales. The traditional hypothesis of cetacean evolution, first proposed by Van Valen in 1966, was that whales were related to the mesonychians, an extinct order of carnivorous ungulates (hoofed animals) that resembled wolves with hooves and were a sister group of the artiodactyls (even-toed ungulates). This hypothesis was proposed due to similarities between the unusual triangular teeth of the mesonychians and those of early whales. However, molecular phylogeny data indicates that whales are very closely related to the artiodactyls, with hippopotamuses as their closest living relative. Because of this observation, cetaceans and hippopotamuses are placed in the same suborder, Whippomorpha. Cetartiodactyla (formed from the words Cetacea and Artiodactyla) is a proposed name for an order that includes both cetaceans and artiodactyls. However, the earliest anthracotheres, the ancestors of hippos, do not appear in the fossil record until the Middle Eocene, millions of years after Pakicetus, whereas the first known whale ancestor appeared during the Early Eocene; this difference in timing implies that the two groups diverged well before the Eocene. Molecular analysis identifies artiodactyls as being very closely related to cetaceans, so mesonychians are probably an offshoot from Artiodactyla, and cetaceans did not derive directly from mesonychians, but the two groups may share a common ancestor.
The molecular data are supported by the discovery of Pakicetus, the earliest archaeocete. The skeletons of Pakicetus show that whales did not derive directly from mesonychians. Instead, they are artiodactyls that began to take to the water soon after artiodactyls split from mesonychians. Archaeocetes retained aspects of their mesonychian ancestry (such as the triangular teeth) which modern artiodactyls, and modern whales, have lost. The earliest ancestors of all hoofed mammals were probably at least partly carnivorous or scavengers, and today's artiodactyls and perissodactyls became herbivores later in their evolution. Whales, however, retained their carnivorous diet because prey was more available and they needed higher caloric content in order to live as marine endotherms (warm-blooded). Mesonychians also became specialized carnivores, but this was likely a disadvantage because large prey was uncommon. This may be why they were out-competed by better-adapted animals like the hyaenodontids and later Carnivora.
Indohyus
Indohyus was a small chevrotain-like animal that lived about 48 million years ago in what is now Kashmir. It belongs to the artiodactyl family Raoellidae, which is believed to be the closest sister group of Cetacea. Indohyus is identified as an artiodactyl because it has two trochlea hinges, a trait unique to artiodactyls. Approximately the size of a raccoon or domestic cat, this omnivorous creature shared some traits of modern whales, most notably the involucrum, a bone growth pattern which is the diagnostic characteristic of any cetacean; this is not found in any other species. It also showed signs of adaptations to aquatic life, including dense limb bones that reduce buoyancy so that they could stay underwater, which are similar to the adaptations found in modern aquatic mammals such as the hippopotamus. This suggests a similar survival strategy to the African chevrotain or water chevrotain which, when threatened by a bird of prey, dives into water and hides beneath the surface for up to four minutes.
The first fossils of the Indohyus were unearthed by Indian geologist A. Ranga Rao. He discovered a few teeth and a jawbone amongst rocks that he had collected. After his death, his widow Leelavathi Rao donated the rocks to professor Hans Thewissen. His technician accidentally broke open a couple of the donated rocks and discovered additional Indohyus fossils.
Pakicetidae
The pakicetids were digitigrade hoofed mammals that are thought to be the earliest known cetaceans, with Indohyus being the closest sister group. They lived in the early Eocene, around 50 million years ago. Their fossils were first discovered in North Pakistan in 1979, located at a river not far from the shores of the former Tethys Sea. After the initial discovery, more fossils were found, mainly in the early Eocene fluvial deposits in northern Pakistan and northwestern India. Based on this discovery, pakicetids most likely lived in an arid environment with ephemeral streams and moderately developed floodplains millions of years ago. By using stable oxygen isotopes analysis, they were shown to drink fresh water, implying that they lived around freshwater bodies. Their diet probably included land animals that approached water for drinking or some freshwater aquatic organisms that lived in the river. The elongated cervical vertebrae and the four, fused sacral vertebrae are consistent with artiodactyls, making Pakicetus one of the earliest fossils to be recovered from the period following the Cetacea/Artiodactyla divergence event.
Pakicetids are classified as cetaceans mainly due to the structure of the auditory bulla (ear bone), which is formed only from the ectotympanic bone. The shape of the ear region in pakicetids is highly unusual and the skull is cetacean-like, although a blowhole is still absent at this stage. The jawbone of pakicetids also lacks the enlarged space (mandibular foramen) that is filled with fat or oil, which is used in receiving underwater sound in modern cetaceans. They have dorsal orbits (eye sockets facing up), which are similar to crocodiles. This eye placement helps submerged predators observe potential prey above the water. According to a 2009 study, the teeth of pakicetids also resemble the teeth of fossil whales, being less like a dog's incisors, and having serrated triangular teeth, which is another link to more modern cetaceans.
It was initially thought that the ears of pakicetids were adapted for underwater hearing, but, as would be expected from the anatomy of the rest of this creature, the ears of pakicetids are specialized for hearing on land. However, pakicetids were able to listen underwater by using enhanced bone conduction, rather than depending on the tympanic membrane like other land mammals. This method of hearing did not give directional hearing underwater.
Pakicetids have long thin legs, with relatively short hands and feet which suggest that they were poor swimmers. To compensate for that, their bones are unusually thick (osteosclerotic), which is probably an adaptation to make the animal heavier to counteract the buoyancy of the water. According to a 2001 morphological analysis by Thewissen et al., pakicetids display no aquatic skeletal adaptation; instead they display adaptations for running and jumping. Hence pakicetids were most likely aquatic waders.
Ambulocetidae
Ambulocetus, which lived about 49 million years ago, was discovered in Pakistan in 1994. They were vaguely crocodile-like mammals, possessing large brevirostrine jaws. In the Eocene, ambulocetids inhabited the bays and estuaries of the Tethys Sea in northern Pakistan. The fossils of ambulocetids are always found in near-shore shallow marine deposits associated with abundant marine plant fossils and littoral mollusks. Although they are found only in marine deposits, their oxygen isotope values indicate that they consumed a range of water with different degrees of salinity, with some specimens having no evidence of sea water consumption and others that did not ingest fresh water at the time when their teeth were fossilized. It is clear that ambulocetids tolerated a wide range of salt concentrations. Hence, ambulocetids represent a transition phase of cetacean ancestors between fresh water and marine habitat.
The mandibular foramen in ambulocetids had increased in size, which indicates that a fat pad was likely to be housed in the lower jaw. In modern toothed whales, this fat pad in the mandibular foramen extends posteriorly to the middle ear. This allows sounds to be received in the lower jaw, and then transmitted through the fat pad to the middle ear. Similar to pakicetids, the orbits of ambulocetids are on the top of the skull, but they face more laterally than in pakicetids.
Ambulocetids had relatively long limbs with particularly strong hind legs, and they retained a tail with no sign of a fluke. The hindlimb structure of Ambulocetids shows that their ability to engage in terrestrial locomotion was significantly limited compared to that of contemporary terrestrial mammals, and likely did not come to land at all. The skeletal structures of the knee and ankle indicates that the motion of the hindlimbs was restricted into one plane. This suggests that, on land, propulsion of the hindlimbs was powered by the extension of dorsal muscles. They probably swam by pelvic paddling (a way of swimming which mainly utilizes their hind limbs to generate propulsion in water) and caudal undulation (a way of swimming which uses the undulations of the vertebral column to generate force for movements), as otters, seals and modern cetaceans do. This is an intermediate stage in the evolution of cetacean locomotion, as modern cetaceans swim by caudal oscillation (a way of swimming similar to caudal undulation, but is more energy efficient).
A recent study suggests that ambulocetids were fully aquatic like modern cetaceans, possessing a similar thoracic morphology and being unable to support their weight on land. This suggests that complete abandonment of the land evolved much earlier among cetaceans than previously thought. However the scientists involved in the study cautioned that the study was limited by a lack of information on the exact density of the bone, the location of the centre of mass, and the reliance of false ribs for thoracic support.
Remingtonocetidae
Remingtonocetids lived in the Middle-Eocene in South Asia, about 49 to 43 million years ago. Compared to family Pakicetidae and Ambulocetidae, Remingtonocetidae was a diverse family found in north and central Pakistan and western India. Remingtonocetids were also found in shallow marine deposits, but they were obviously more aquatic than ambulocetidae. This is demonstrated by the recovery of their fossils from a variety of coastal marine environments, including near-shore and lagoonal deposits. According to stable oxygen isotopes analysis, most remingtonocetids did not ingest fresh water, and had hence lost their dependency on fresh water relatively soon after their origin.
The orbits of remingtonocetids faced laterally and were small. This suggests that vision was not an important sense for them. The nasal opening, which eventually becomes the blowhole in modern cetaceans, was located near the tip of the snout. The position of the nasal opening had remained unchanged since pakicetids. One of the notable features in remingtonocetids is that the semicircular canals, which are important for balancing in land mammals, had decreased in size. This reduction in size had closely accompanied the cetacean radiation into marine environments. According to a 2002 study done by Spoor et al., this modification of the semicircular canal system may represent a crucial 'point of no return' event in early cetacean evolution, which excluded a prolonged semi-aquatic phase.
Compared to ambulocetids, remingtonocetids had relatively short limbs. Based on their skeletal remains, remingtonocetids were probably amphibious cetaceans that were well adapted to swimming, and likely to swim by caudal undulation only.
Protocetidae
The protocetids form a diverse and heterogeneous group known from Asia, Europe, Africa, and North America. They lived in the Eocene, approximately 48 to 35 million years ago. The fossil remains of protocetids were uncovered from coastal and lagoonal facies in South Asia; unlike previous cetacean families, their fossils uncovered from Africa (e.g., Phiomicetus, Protocetus, Aegyptocetus, Togocetus) and North America (e.g., Georgiacetus) also include open marine forms. They were probably amphibious, but more aquatic compared to remingtonocetids. Protocetids were the first cetaceans to leave the Indian subcontinent and disperse to all shallow subtropical oceans of the world.
There were many genera among the family Protocetidae. There were different degrees of aquatic adaptations in this group, with some able to support their weight on land, and others not. The discovery of the southeastern Pacific Peregocetus indicates they crossed the Atlantic and achieved circumstances-equatorial distribution by 40 mya. Their amphibious nature is supported by the discovery of a pregnant Maiacetus, in which the fossilised fetus was positioned for a head-first delivery, suggesting that Maiacetus gave birth on land. If they gave birth in the water, the fetus would be positioned for a tail-first delivery to avoid drowning during birth.
Unlike remingtonocetids and ambulocetids, protocetids have large orbits which are oriented laterally. Increasingly lateral-facing eyes might be used to observe underwater prey, and are similar to the eyes of modern cetaceans. Furthermore, the nasal openings were large and were halfway up the snout. The great variety of teeth suggests diverse feeding modes in protocetids.
In both remingtonocetids and protocetids, the size of the mandibular foramen had increased. The large mandibular foramen indicates that the mandibular fat pad was present. However the air-filled sinuses that are present in modern cetaceans, which function to isolate the ear acoustically to enable better underwater hearing, were still not present. The external auditory meatus (ear canal), which is absent in modern cetaceans, was also present. Hence, the method of sound transmission that were present in them combines aspects of pakicetids and modern odontocetes (toothed whales). At this intermediate stage of hearing development, the transmission of airborne sound was poor due to the modifications of the ear for underwater hearing while directional underwater hearing was also poor compared to modern cetaceans.
Some protocetids had short, wide fore- and hindlimbs that were likely to have been used in swimming, but the limbs gave a slow and cumbersome locomotion on land. It is possible that some protocetids had flukes. However, it is clear that they were adapted even further to an aquatic life-style. In Rodhocetus, for example, the sacrum (a bone that, in land-mammals, is a fusion of five vertebrae that connects the pelvis with the rest of the vertebral column) was divided into loose vertebrae. However, the pelvis was still connected to one of the sacral vertebrae. The ungulate ancestry of these archaeocetes is still underlined by characteristics like the presence of hooves at the ends of the toes in Rodhocetus.
The foot structure of Rodhocetus shows that protocetids were predominantly aquatic. A 2001 study done by Gingerich et al. hypothesized that Rodhocetus locomoted in the oceanic environment similarly to how ambulocetids pelvic paddling, which was supplemented by caudal undulation. Terrestrial locomotion of Rodhocetus was very limited due to their hindlimb structure. It is thought that they moved in a way similar to how eared seals move on land, by rotating their hind flippers forward and underneath their body.
Basilosauridae
Basilosaurids and dorudontines lived together in the late Eocene around 41 to 33.9 million years ago, and are the oldest known obligate aquatic cetaceans. They were fully recognizable whales which lived entirely in the ocean. This is supported by their fossils usually found in deposits indicative of fully marine environments, lacking any freshwater influx. They were probably distributed throughout the tropical and subtropical seas of the world. Basilosaurids are commonly found in association with dorudontines, and were closely related to one another. The fossilised stomach contents in one basilosaurid indicates that it ate fish.
Although they look very much like modern cetaceans, basilosaurids lacked the 'melon organ' that allows toothed whales to use echolocation. They had small brains; this suggests they were solitary and did not have the complex social structures of some modern cetaceans. The mandibular foramen of basilosaurids covered the entire depth of the lower jaw as in modern cetaceans. Their orbits faced laterally, and the nasal opening had moved even higher up the snout, closer to the position of the blowhole in modern cetaceans. Furthermore, their ear structures were functionally modern, with the insertion of air-filled sinuses between ear and skull. Unlike modern cetaceans, basilosaurids retained a large external auditory meatus.
Both basilosaurids have skeletons that are immediately recognizable as cetaceans. A basilosaurid was as big as the larger modern whales, with genera like Basilosaurus reaching lengths of up to long; dorudontines were smaller, with genera like Dorudon reaching about long. The smallest basilosaurid whale is Tutcetus and measures long. The large size of basilosaurids is due to the extreme elongation of their lumbar vertebrae. They had a tail fluke, but their body proportions suggest that they swam by caudal undulation and that the fluke was not used for propulsion. In contrast, dorudontines had a shorter but powerful vertebral column. They too had a fluke and, unlike basilosaurids, they probably swam similarly to modern cetaceans, by using caudal oscillations. The forelimbs of basilosaurids were probably flipper-shaped, and the external hind limbs were tiny and were certainly not involved in locomotion. Their fingers, however, retained the mobile joints of their ambulocetid relatives. The two tiny but well-formed hind legs of basilosaurids were probably used as claspers when mating. The pelvic bones associated with these hind limbs were not connected to the vertebral column as they were in protocetids. Essentially, any sacral vertebrae can no longer be clearly distinguished from the other vertebrae.
Both basilosaurids and dorudontines are relatively closely related to modern cetaceans, which belong to parvorders Odontoceti and Mysticeti. However, according to a 1994 study done by Fordyce and Barnes, the large size and elongated vertebral body of basilosaurids preclude them from being ancestral to extant forms. As for dorudontines, there are some species within the family that do not have elongated vertebral bodies, which might be the immediate ancestors of Odontoceti and Mysticeti. The other basilosaurids became extinct.
Evolution of modern cetaceans
Baleen whales
All modern baleen whales or mysticetes are filter-feeders which have baleen in place of teeth, though the exact means by which baleen is used differs among species (gulp-feeding within balaenopterids, skim-feeding within balaenids, and bottom plowing within eschrichtiids). The first members of both groups appeared during the middle Miocene. Filter feeding is very beneficial as it allows baleen whales to efficiently gain huge energy resources, which makes the large body size in modern varieties possible. The development of filter feeding may have been a result of worldwide environmental change and physical changes in the oceans. A large-scale change in ocean current and temperature could have contributed to the radiation of modern mysticetes. The earlier varieties of baleen whales, or "archaeomysticetes", such as Janjucetus and Mammalodon had very little baleen and relied mainly on their teeth.
There is also evidence of a genetic component of the evolution of toothless whales. Multiple mutations have been identified in genes related to the production of enamel in modern baleen whales. These are primarily insertion/deletion mutations that result in premature stop codons. It is hypothesized that these mutations occurred in cetaceans already possessing preliminary baleen structures, leading to the pseudogenization of a "genetic toolkit" for enamel production. Recent research has also indicated that the development of baleen and the loss of enamel-capped teeth both occurred once, and both occurred on the mysticete stem branch.
Generally it is speculated the four modern mysticete families have separate origins among the cetotheres. Modern baleen whales, Balaenopteridae (rorquals and humpback whale, Megaptera novaengliae), Balaenidae (right whales), Eschrichtiidae (gray whale, Eschrictius robustus), and Neobalaenidae (pygmy right whale, Caperea marginata) all have derived characteristics presently unknown in any cetothere and vice versa (such as a sagittal crest). Mysticetes are also known for their gigantism, as baleen whales are among the largest organisms to ever have lived; they reach lengths greater than 20 m and weigh more than 100,000 kg. This gigantism is directly related to their feeding mechanism – mysticete size has been found to be dependent on the amount of baleen a mysticete can use to filter its prey. Additionally, size is a positively selected trait that gives mysticetes a boost in fitness. Mysticete populations will therefore slowly become even more gigantic as whales with larger amounts of baleen are selected.
Toothed whales
The adaptation of echolocation occurred when toothed whales (Odontoceti) split apart from baleen whales, and distinguishes modern toothed whales from fully aquatic archaeocetes. This happened around 34 million years ago in a second cetacean radiation. Modern toothed whales do not rely on their sense of sight, but rather on their sonar to hunt prey. Echolocation also allowed toothed whales to dive deeper in search of food, with light no longer necessary for navigation, which opened up new food sources. Toothed whales echolocate by creating a series of clicks emitted at various frequencies. Sound pulses are emitted, reflected off objects, and retrieved through the lower jaw. Skulls of Squalodon show evidence for the first hypothesized appearance of echolocation. Squalodon lived from the early to middle Oligocene to the middle Miocene, around 33–14 million years ago. Squalodon featured several commonalities with modern toothed whales: the cranium was well compressed (to make room for the melon, a part of the nose), the rostrum telescoped outward into a beak, a characteristic of the modern toothed whales that gave Squalodon an appearance similar to them. However, it is thought unlikely that squalodontids are direct ancestors of modern toothed whales.
The first oceanic dolphins such as kentriodonts, evolved in the late Oligocene and diversified greatly during the mid-Miocene. The first fossil cetaceans near shallow seas (where porpoises inhabit) were found around the North Pacific; species like Semirostrum were found along California (in what were then estuaries). These animals spread to the European coasts and Southern Hemisphere only much later, during the Pliocene. The earliest known ancestor of arctic whales is Denebola brachycephala from the late Miocene around 9–10 million years ago. A single fossil from Baja California indicates the family once inhabited warmer waters.
Ancient sperm whales differ from modern sperm whales in tooth count and the shape of the face and jaws. For example, Scaldicetus had a tapered rostrum. Genera from the Oligocene and Miocene had teeth in their upper jaws. These anatomical differences suggest that these ancient species may not have necessarily been deep-sea squid hunters like the modern sperm whale, but that some genera mainly ate fish. Contrary to modern sperm whales, most ancient sperm whales were built to hunt whales. Livyatan had a short and wide rostrum measuring across, which gave the whale the ability to inflict major damage on large struggling prey, such as other early whales. Species like these are collectively known as killer sperm whales or macroraptorial sperm whales.
Beaked whales consist of over 20 genera. Earlier variety were probably preyed upon by killer sperm whales and large sharks such as megalodon. In 2008, a large number of fossil ziphiids were discovered off the coast of South Africa, confirming the remaining ziphiid species might just be a remnant of a higher diversity that has since gone extinct. After studying numerous fossil skulls, researchers discovered the absence of functional maxillary teeth in all South African ziphiids, which is evidence that suction feeding had already developed in several beaked whale lineages during the Miocene. Extinct ziphiids also had robust skulls, suggesting that tusks were used for male-male interactions.
Skeletal evolution
Modern cetaceans have internal, rudimentary hind limbs, such as reduced femurs, fibulas, and tibias, and a pelvic girdle. Indohyus has a thickened ectotympanic internal lip of the ear bone. This feature compares directly to that of modern cetaceans. Another similar feature was the composition of the teeth, which contained mostly calcium phosphate which is needed for eating and drinking by aquatic animals, though, unlike modern day toothed whales, they had a heterodont (more than one tooth morphology) dentition as opposed to a homodont (one tooth morphology present) dentition. Although they somewhat resembled a wolf, the fossils of pakicetids showed the eye sockets were much closer to the top of their head than that of other terrestrial mammals, but similar to the structure of the eyes in cetaceans. Their transition from land to water led to reshaping of the skull and food processing equipment because the eating habits were changing. The change in position of the eyes and limb bones is associated with the pakicetids becoming waders. The ambulocetids also began to develop long snouts, which is seen in current cetaceans. Their limbs (and hypothesized movement) were very similar to otters.
Limblessness in cetaceans does not represent a regression of fully formed limbs nor the absence of limb bud initiation, but rather arrest of limb bud development. Limb buds develop normally in cetacean embryos. Limb buds progress to the condensation phase of early skeletogenesis, where nerves grow into the limb bud and the apical ectodermal ridge (AER), a structure that ensures proper limb development, appears functional. Occasionally, the genes that code for longer extremities cause a modern whale to develop miniature legs (atavism).
Pakicetus had a pelvic bone most similar to that of terrestrial mammals. In later species, such as Basilosaurus, the pelvic bone, no longer attached to the vertebrae and the ilium, was reduced. Certain genes are believed to be responsible for the changes that occurred to the cetacean pelvic structure, such as BMP7, PBX1, PBX2, PRRX1, and PRRX2. The pelvic girdle in modern cetaceans were once thought to be vestigial structures that served no purpose at all. The pelvic girdle in male cetaceans is different in size compared to females, and the size is thought to be a result of sexual dimorphism. The pelvic bones of modern male cetaceans are more massive, longer, and larger than those of females. Due to the sexual dimorphism displayed, they were most likely involved in supporting male genitalia that remain hidden behind abdominal walls until sexual reproduction occurs.
Early archaeocetes such as Pakicetus had the nasal openings at the end of the snout, but in later species such as Rodhocetus, the openings had begun to drift toward the top of the skull. This is known as nasal drift. The nostrils of modern cetaceans have become modified into blowholes that allow them to break to the surface, inhale, and submerge with convenience. The ears began to move inward as well, and, in the case of Basilosaurus, the middle ears began to receive vibrations from the lower jaw. Today's modern toothed whales use their melon organ, a pad of fat, for echolocation.
Radiation events
There are three major radiation events that mark diversification and speciation in the evolutionary history of Cetacea. The first occurred around the middle Eocene (40 Mya) when these early cetaceans abandoned riverine and shallow coastal habitats, setting the scene for Protocetidae – the first fully marine cetacean. With the oceans and its nutrients at their disposal, rapidly diversifying protocetids were also responsible for the first major geographic expansion, dispersing throughout North Africa, Europe, and North America. The second of three major radiation events occurred near the start of the Oligocene (~34 Mya) when Neoceti diverged from Basilosauridae. This radiation event concurrently occurs with the breakup of Gondwana and the opening of the Southern Ocean, wildly changing ocean ecosystems, productivity, and temperature gradients. The timing of this second radiation event is not coincidental, as the following diversification of cetaceans was likely due to new ecological opportunities the change in oceans gave them. The final major radiation event, occurring throughout the middle Miocene and into the Pliocene (12 Mya to 2 Mya), was not due to a specific event but is associated with widespread generic expansion of odontocetes and mysticetes. Some modern genera of cetaceans began to emerge, including Balaenoptera, a genus of rorquals that includes the blue whale. Delphinidae, ocean dolphins, also arose during this radiation event in the late Miocene.
Ongoing evolution
Culture
Culture is group-specific behavior transferred by social learning. Tool use to aid with foraging is one example. Whether or not a dolphin uses a tool affects its eating behavior, which causes differences in diet. Also, using a tool allows a new niche and new prey to open up for that particular dolphin. Due to these differences, fitness levels change within the dolphins of a population, which further causes evolution to occur in the long run. Culture and social networks have played a large role in the evolution of modern cetaceans, as concluded in studies showing dolphins preferring mates with the same socially learned behaviors, and humpback whales using songs between breeding areas. For dolphins particularly, the largest non-genetic effects on their evolution are due to culture and social structure.
Based on a 2014 study, the population of Indo-Pacific bottlenose dolphins (Tursiops sp.) around Shark Bay of Western Australia can be divided into spongers and nonspongers. Spongers put sea sponges on their snouts as protection against abrasions from sharp objects, stingray barbs, or toxic organisms. The sponges also help the dolphins target fish without swim bladders, since echolocation cannot detect these fish easily against a complex background. Spongers also specifically forage in deep channels, but nonspongers are found foraging in both deep and shallow channels. This foraging behavior is mainly passed on from mother to child. Therefore, since this is a group behavior being passed down by social learning, this tool use is considered a cultural trait.
Researchers in a 2014 study in Shark Bay found the fatty acid analyses between the West and East Gulf populations to differ, which is due to the two areas having different food sources. However, when comparing data from within the West Gulf, the spongers vs. the nonspongers in the deep channels had very different fatty acid results even though they are in the same habitat. Nonspongers from deep and shallow channels had similar data. This suggests that sponging was the cause of the different data and not the deep vs. shallow channels. Sponging opened up a new niche for the dolphins and allowed them access to new prey, which caused long-term dietary changes. By producing different food sources within a population, there is less intrapopulation competition for resources, showing character displacement. As a result, the carrying capacity increases since the entire population does not depend on one food source. The fitness levels within the population also change, thus allowing this culture to evolve.
Social structure
Social structure forms groups with individuals that interact with one another, and this allows for cultural traits to emerge, exchange, and evolve. This relationship is especially seen in the bottlenose dolphin populations in southwestern Australia, which have been known to beg for food from fishermen. This begging behavior was spread through the population due to individual (dolphins spending time around boats) and social (dolphins spending time with other dolphins who express begging behavior) learning.
Culture can, however, impact social structure by causing behavior matching and assortive mating. Individuals within a certain culture are more likely to mate with individuals using the same behaviors rather than a random individual, thus influencing social groups and structure. For example, the spongers of Shark Bay preferentially stick with other spongers. Also, some bottlenose dolphins in Moreton Bay, Australia followed prawn trawlers to feed on their debris, while other dolphins in the same population did not. The dolphins preferentially associated with individuals with same behavior even though they all lived in the same habitat. Later on, prawn trawlers were no longer present, and the dolphins integrated into one social network after a couple of years.
Social networks can still affect and cause evolution on their own by impending fitness differences on individuals. According to a 2012 study, male calves had a lower survival rate if they had stronger bonds with juvenile males. However, when other age and sex classes were tested, their survival rate did not significantly change. This suggests that juvenile males impose a social stress on their younger counterparts. In fact, it has been documented that juvenile males commonly perform acts of aggression, dominance, and intimidation against the male calves. According to a 2010 study, certain populations of Shark Bay dolphins had varying levels of fitness and calf success. This is either due to social learning (whether or not the mother passed on her knowledge of reproductive ability to the calves), or due to the strong association between mother dolphins in the population; by sticking in a group, an individual mother does not need to be as vigilant all the time for predators.
Genetic studies conducted on Clymene dolphins (Stenella clymene) focused on their natural histories, and the results show that the origin of the species was actually an outcome of hybrid speciation. Hybridization between spinner dolphins (Stenella longirostris) and striped dolphins (Stenella coeruleoalba) in the North Atlantic was caused by constant habitat sharing of the two species. Relationships between these three species had been speculated according to notable resemblances between anatomies of the Clymene and the spinner dolphins, resulting in the former being regarded as subspecies of the latter until 1981, and the possibility of the Clymene dolphin as a hybrid between the spinner and the striped dolphins have come to question based on anatomical and behavioral similarities between these two species.
Environmental factors
Genome sequences done in 2013 revealed that the Yangtze river dolphin, or "baiji" (Lipotes vexillifer), lacks single nucleotide polymorphisms in their genome. After reconstructing the history of the baiji genome for this dolphin species, researchers found that the major decrease in genetic diversity occurred most likely due to a bottleneck event during the last deglaciation event. During this time period, sea levels were rising while global temperatures were increasing. Other historical climate events can be correlated and matched with the genome history of the Yangtze river dolphin as well. This shows how global and local climate change can drastically affect a genome, leading to changes in fitness, survival, and evolution of a species.
The European population of common dolphins (Delphinus delphis) in the Mediterranean have differentiated into two types: eastern and western. According to a 2012 study, this seems to be due to a recent bottleneck as well, which drastically decreased the size of the eastern Mediterranean population. Also, the lack of population structure between the western and eastern regions seems contradictory of the distinct population structures between other regions of dolphins. Even though the dolphins in the Mediterranean area had no physical barrier between their regions, they still differentiated into two types due to ecology and biology. Therefore, the differences between the eastern and western dolphins most likely stems from highly specialized niche choice rather than just physical barriers. Through this, environment plays a large role in the differentiation and evolution of this dolphin species.
The divergence and speciation within bottlenose dolphins has been largely due to climate and environmental changes over history. According to research, the divisions within the genus correlate with periods of rapid climate change. For example, the changing temperatures could cause the coast landscape to change, niches to empty up, and opportunities for separation to appear. In the Northeast Atlantic, specifically, genetic evidence suggests that the bottlenose dolphins have differentiated into coastal and pelagic types. Divergence seems most likely due to a founding event where a large group separated. Following this event, the separate groups adapted accordingly and formed their own niche specializations and social structures. These differences caused the two groups to diverge and to remain separated.
Two endemic, distinctive types of short-finned pilot whale, Tappanaga (or Shiogondou) the larger, northern type and Magondou the smaller, southern type, can be found along the Japanese archipelago where distributions of these two types mostly do not overlap by the oceanic front border around the easternmost point of Honshu. It is thought that the local extinction of long-finned pilot whales in the North Pacific in the 12th century could have triggered the appearance of Tappanaga, causing short-finned pilot whales to colonize the colder ranges of the long-finned variant. Whales with similar characteristics to the Tappanaga can be found along Vancouver Island and northern US coasts as well.
| Biology and health sciences | Cetaceans | Animals |
875176 | https://en.wikipedia.org/wiki/Colemanite | Colemanite | Colemanite (Ca2B6O11·5H2O) or (CaB3O4(OH)3·H2O) is a borate mineral found in evaporite deposits of alkaline lacustrine environments. Colemanite is a secondary mineral that forms by alteration of borax and ulexite.
It was first described in 1884 for an occurrence near Furnace Creek in Death Valley and was named after William Tell Coleman (1824–1893), owner of the mine "Harmony Borax Works" where it was first found. At the time, Coleman had alternatively proposed the name "smithite" instead after his business associate Francis Marion Smith.
Uses
Colemanite is an important ore of boron, and was the most important boron ore until the discovery of kernite in 1926. It has many industrial uses, like the manufacturing of heat resistant glass.
Occurrence
About 40% of the world's known colemanite reserves are at the Emet mine in western Turkey. Other important sources in Turkey are found at Bigadiç and Kestelek.
| Physical sciences | Minerals | Earth science |
875205 | https://en.wikipedia.org/wiki/Cuprite | Cuprite | Cuprite is an oxide mineral composed of copper(I) oxide Cu2O, and is a minor ore of copper.
Its dark crystals with red internal reflections are in the isometric system hexoctahedral class, appearing as cubic, octahedral, or dodecahedral forms, or in combinations. Penetration twins frequently occur. In spite of its nice color, it is rarely used for jewelry because of its low Mohs hardness of 3.5 to 4. It has a relatively high specific gravity of 6.1, imperfect cleavage and is brittle to conchoidal fracture. The luster is sub-metallic to brilliant adamantine. The "chalcotrichite" (from , "plush copper ore") variety typically shows greatly elongated (parallel to [001]) capillary or needle like crystals forms.
It is a secondary mineral which forms in the oxidized zone of copper sulfide deposits. It frequently occurs in association with native copper, azurite, chrysocolla, malachite, tenorite and a variety of iron oxide minerals. It is known as ruby copper due to its distinctive red color.
Cuprite was first described by Wilhelm Karl Ritter von Haidinger in 1845 and the name derives from the Latin cuprum for its copper content.
Cuprite is found in the Ural Mountains, Altai Mountains, and Sardinia, and in more isolated locations in Cornwall, France, Arizona, Chile, Bolivia, and Namibia.
As a gemstone
Though almost all crystals of cuprite are far too small to yield faceted gemstones, one unique deposit from Onganja in Seeis, Namibia, which was discovered in the 1970s, has produced crystals which were both large and gem quality. Virtually every faceted stone over one carat (0.2 g) in weight is from this single deposit, which has long since been mined out. The number of faceted gems over two carats (0.4 g) is difficult to estimate, but according to Joel Arem, one-time curator for the Smithsonian National Gem and Mineral Collection in Washington, D.C., faceted cuprite of any size is considered one of the most collectible and spectacular gems in existence, with its deep garnet coloring and higher brilliance than a diamond. Only the gem's soft nature prevents it from being among the most valuable jewelry stones.
| Physical sciences | Minerals | Earth science |
875336 | https://en.wikipedia.org/wiki/Realgar | Realgar | Realgar ( ), also known as arsenic blende, ruby sulphur or ruby of arsenic, is an arsenic sulfide mineral with the chemical formula α-. It is a soft, sectile mineral occurring in monoclinic crystals, or in granular, compact, or powdery form, often in association with the related mineral, orpiment (). It is orange-red in color, melts at 320 °C, and burns with a bluish flame releasing fumes of arsenic and sulfur. Realgar is soft with a Mohs hardness of 1.5 to 2 and has a specific gravity of 3.5. Its streak is orange colored. It is trimorphous with pararealgar and bonazziite.
Etymology
Its name comes from the Arabic rahj al-ġār ( , "powder of the mine"), via Medieval Latin, and its earliest record in English is in the 1390s.
Uses
Realgar is a minor ore of arsenic extracted in China, Peru, and the Philippines.
Historical uses
Realgar was used by firework manufacturers to create the color white in fireworks prior to the availability of powdered metals such as aluminium, magnesium and titanium. It is still used in combination with potassium chlorate to make a contact explosive known as "red explosive" for some types of torpedoes and other novelty exploding fireworks branded as 'cracker balls', as well in the cores of some types of crackling stars.
Realgar is toxic. It was sometimes used to kill weeds, insects, and rodents, even though more effective arsenic-based anti-pest agents are available such as cacodylic acid, , an organoarsenic compound used as herbicide.
Realgar was commonly used in leather manufacturing to remove hair from animal pelts. Because it is a known carcinogen and an arsenic poison, and because substitutes are available, it is rarely used today for this purpose.
The ancient Greeks, who called realgar (), understood that it was poisonous. From this, realgar has also historically been known in English as sandarac.
Realgar was also used by Ancient Greek apothecaries to make a medicine known as "bull's blood". The Greek physician Nicander described a death by "bull's blood", which matches the known effects of arsenic poisoning. Bull's blood is the poison that is said to have been used by Themistocles and Midas for suicide.
The Chinese name for realgar is (Mandarin ), literally 'masculine yellow', as opposed to orpiment which is 'feminine yellow'.
Realgar was, along with orpiment, traded in the Roman Empire and was used as a red paint pigment. Early occurrences of realgar as a red paint pigment are known for works of art from China, India, Central Asia, and Egypt. It was used in Venetian fine-art painting during the Renaissance era, though rarely elsewhere in Europe, a use which died out by the 18th century. It was also used as medicine. Other traditional uses include manufacturing lead shot, printing, and dyeing calico cloth. It was used to poison rats in medieval Spain and in 16th century England.
Occurrence
Realgar most commonly occurs as a low-temperature hydrothermal vein mineral associated with other arsenic and antimony minerals. It also occurs as volcanic sublimations and in hot spring deposits. It occurs in association with orpiment, arsenolite, calcite and barite.
It is found with lead, silver and gold ores in Hungary, Bohemia and Saxony. In the US it occurs notably in Mercur, Utah; Manhattan, Nevada; and in the geyser deposits of Yellowstone National Park.
After a long period of exposure to light, realgar changes form to a yellow powder known as pararealgar (β-). It was once thought that this powder was the yellow sulfide orpiment, but is a distinct chemical compound.
Gallery
| Physical sciences | Minerals | Earth science |
875354 | https://en.wikipedia.org/wiki/Phyllite | Phyllite | Phyllite ( ) is a type of foliated metamorphic rock formed from slate that is further metamorphosed so that very fine grained white mica achieves a preferred orientation. It is primarily composed of quartz, sericite mica, and chlorite.
Phyllite has fine-grained mica flakes, whereas slate has extremely fine mica flakes, and schist has large mica flakes, all mica flakes of which have achieved a preferred orientation. Among foliated metamorphic rocks, it represents a gradation in the degree of metamorphism between slate and schist.
The minute crystals of graphite, sericite, or chlorite, or the translucent fine-grained white mica, impart a silky, sometimes golden sheen to the surfaces of cleavage, called "phyllitic luster".
The word comes from the Greek phyllon, meaning "leaf".
The protolith (or parent rock) for phyllite is shale or pelite; or slate, which in turn came from a shale protolith. Its constituent platy minerals are larger than those in slate but are not visible with the naked eye. Phyllites are said to have a texture called "phyllitic sheen," and are usually classified as having formed through low-grade metamorphic conditions through regional metamorphism metamorphic facies.
Phyllite has good fissility (a tendency to split into sheets). Phyllites are usually black to gray or light greenish gray in color. The foliation is commonly crinkled or wavy in appearance. Phyllites are mostly used in decorative aggregates, interior decors, building stones, facing stones, garden decoration and curbing. Cemetery markers, commemorative tablets, creating artworks and writing slates are some of its commercial uses.
Phyllite is commonly found in the Dalradian metasediments of northwest Arran. In north Cornwall, there are Tredorn phyllites and Woolgarden phyllites.
Carolina "slate" is often volcanic phyllite. A type of Carolina slate, Duke stone, is a dacitic phyllite that is fractured and colored with iron oxide.
| Physical sciences | Metamorphic rocks | Earth science |
875372 | https://en.wikipedia.org/wiki/Pyrolusite | Pyrolusite | Pyrolusite is a mineral consisting essentially of manganese dioxide (MnO2) and is important as an ore of manganese. It is a black, amorphous appearing mineral, often with a granular, fibrous, or columnar structure, sometimes forming reniform crusts. It has a metallic luster, a black or bluish-black streak, and readily soils the fingers. The specific gravity is about 4.8. Its name is from the Greek for fire and to wash, in reference to its use as a way to remove tints from glass.
Occurrence
Pyrolusite and romanechite are among the most common manganese minerals.
Pyrolusite occurs associated with manganite, hollandite, hausmannite, braunite, cryptomelane, chalcophanite, goethite, and hematite under oxidizing conditions in hydrothermal deposits. It also occurs in bogs and often results from alteration of manganite.
Use
The metal is obtained by reduction of the oxide with sodium, magnesium, aluminium, or by electrolysis. Pyrolusite is extensively used for the manufacture of spiegeleisen and ferromanganese and of various alloys such as manganese-bronze. As an oxidizing agent it is used in the preparation of chlorine; indeed, chlorine gas itself was first described by Karl Scheele in 1774 from the reaction products of pyrolusite and hydrochloric acid. Natural pyrolusite has been used in batteries, but high-quality batteries require synthetic products. Pyrolusite is also used to prepare disinfectants (permanganates) and for decolorizing glass. When mixed with molten glass it oxidizes the ferrous iron to ferric iron, and so discharges the green and brown tints (making it classically useful to glassmakers as a decolorizer). As a coloring material, it is used in calico printing and dyeing; for imparting violet, amber, and black colors to glass, pottery, and bricks; and in the manufacture of green and violet paints.
Dendritic manganese oxides
Black, manganese oxides with a dendritic crystal habit often found on fracture or rock surfaces are often assumed to be pyrolusite although careful analyses of numerous examples of these dendrites has shown that none of them are, in fact, pyrolusite. Instead, they are other forms of manganese oxide.
History
Some of the most famous early cave paintings in Europe were executed by means of manganese dioxide. Blocks of pyrolusite are found often at Neanderthal sites. It may have been kept as a pigment for cave paintings, but it has also been suggested that it was powdered and mixed with tinder fungus for lighting fires. Manganese dioxide, in the form of umber, was one of the earliest natural substances used by human ancestors. It was used as a pigment at least from the Middle Paleolithic. It may have been also used by the Neanderthals in fire-making.
The ancient Greeks had a term μάγνης or Μάγνης λίθος ("Magnes lithos") meaning stone of the area called Μαγνησία (Magnesia), referring to Magnesia in Thessaly or to areas in Asia Minor with that name. Two minerals are called μάγνης, namely lodestone and pyrolusite (manganese dioxide).(not to be confused with the current mineral called Lodestone, which does not contain Manganese. Lodestone is a naturally magnetized form of the iron mineral Magnetite. FeFe2O4)
The term μαγνησία was used for manganese dioxide. In the sixteenth century it was called "manganesum". It also was called Alabandicus (from the Alabanda region of Asia Minor) and Braunstein. Eventually the name of the element manganese was derived from "manganesum", whereas "magnesia" came to mean the oxide of a different element, magnesium.
| Physical sciences | Minerals | Earth science |
875405 | https://en.wikipedia.org/wiki/Scheelite | Scheelite | Scheelite is a calcium tungstate mineral with the chemical formula CaWO4. It is an important ore of tungsten (wolfram). Scheelite is originally named after Swedish chemist Carl Wilhelm Scheele (1742–1786). Well-formed crystals are sought by collectors and are occasionally fashioned into gemstones when suitably free of flaws. Scheelite has been synthesized using the Czochralski process; the material produced may be used to imitate diamond, as a scintillator, or as a solid-state lasing medium. It was also used in radium paint in the same fashion as was zinc sulphide, and Thomas Edison invented a fluoroscope with a calcium tungstate-coated screen, making the images six times brighter than those with barium platinocyanide; the latter chemical allowed Röntgen to discover X-rays in early November 1895. Note, the semi-precious stone marketed as 'blue scheelite' is actually a rock type consisting mostly of calcite and dolomite, with occasional traces of yellow-orange scheelite.
Properties
Its crystals are in the tetragonal crystal system, appearing as dipyramidal pseudo-octahedra. Colors include golden yellow, brownish green to dark brown, pinkish to reddish gray, orange and colorless. Transparency ranges from translucent to transparent and crystal faces are highly lustrous (vitreous to adamantine). Scheelite possesses distinct cleavage and its fracture may be subconchoidal to uneven. Its specific gravity is high at 5.9–6.1 and its hardness is low at 4.5–5. Aside from pseudo-octahedra, scheelite may be columnar, granular, tabular or massive in habit. Druzes are quite rare and occur almost exclusively at Zinnwald, Czech Republic. Twinning is also commonly observed and crystal faces may be striated. Scheelite streaks white and is brittle.
Gems cut from transparent material are fragile. Scheelite's refractive index (1.918–1.937 uniaxial positive, with a maximum birefringence of 0.016) and dispersion (0.026) are both moderately high. These factors combine to result in scheelite's high lustre and perceptible "fire", approaching that of diamond.
Scheelite fluoresces under shortwave ultraviolet light, the mineral glows a bright sky-blue. The presence of molybdenum trace impurities occasionally results in a green glow. Fluorescence of scheelite, sometimes associated with native gold, is used by geologists in the search for gold deposits.
Occurrence
Scheelite occurs in contact metamorphic skarns; in high-temperature hydrothermal veins and greisen; less commonly in granite pegmatites. Temperature and pressure of formation is between and from . Typical mineral association includes cassiterite, wolframite, topaz, fluorite, apatite, tourmaline, quartz, grossular–andradite, diopside, vesuvianite and tremolite.
Scheelite usually occurs in tin-bearing veins; and is sometimes found in association with gold. Fine crystals have been obtained from Caldbeck Fells in Cumbria, Zinnwald/Cínovec and Elbogen in Bohemia, Guttannen in Switzerland, the Giant Mountains in Silesia, Dragoon Mountains in Arizona and elsewhere. At Trumbull in Connecticut and Kimpu-san in Japan large crystals of scheelite completely altered to wolframite have been found: those from Japan have been called “reinite.” It was mined until 1990 at King Island, Australia, Glenorchy in Central Otago and Macraes Flat in North Otago and also at The Golden Bar mine at Dead Horse Creek during World War I in Nelson, New Zealand. There is a high concentration of Scheelite in Northeast of Brazil, mainly in the Currais Novos mine in Rio Grande do Norte State. One of the world's largest Scheelite mining companies is in Luoyang, China.
History
Scheelite was first described in 1751 for an occurrence in Mount Bispbergs klack, Säter, Dalarna, Sweden, and named for Carl Wilhelm Scheele (1742–1786). Owing to its unusual heaviness, it had been given the name tungsten by the Swedes, meaning “heavy stone.” The name was later used to describe the metal, while the ore itself was given the name scheelerz or scheelite.
Synthetics
Although it is now uncommon as a diamond imitation (much more convincing products, like cubic zirconia and moissanite have long since superseded it), synthetic scheelite is occasionally offered as natural scheelite, and collectors may thus be fooled into paying highprices for them. Gemologists distinguish natural scheelite from synthetic material mainly by microscopic examination: Natural material is very seldom without internal growth features and inclusions (imperfections), while synthetic material is usually very clean. Distinctly artificial curved striae and clouds of minute gas bubbles may also be observed in synthetic scheelite.
The visible absorption spectrum of scheelite, as seen by a hand-held (direct-vision) spectroscope, may also be of use: most natural stones show a number of faint absorption lines in the yellow region of the spectrum (~585 nm) due to praseodymium and neodymium trace impurities. Conversely, synthetic scheelite is often without such a spectrum. Some synthetics may however be doped with neodymium or other rare-earth elements, but the spectrum produced is unlike that of natural stones.
Applications
Scheelite is widely used in phosphors, particularly in scintillators for X-ray and gamma-ray detection. It is also utilized in fluorescent lighting systems for its ability to convert ultraviolet light into visible light. In some cathode ray tubes (CRTs), calcium tungstate (Scheelite) is used as a phosphorescent screen material.
In popular culture
Scheelite figures in the manga series Dr. Stone, as a precursor to tungsten, and for its fluorescence.
| Physical sciences | Minerals | Earth science |
875418 | https://en.wikipedia.org/wiki/Siderite | Siderite | Siderite is a mineral composed of iron(II) carbonate (FeCO3). Its name comes from the Ancient Greek word (), meaning "iron". A valuable iron ore, it consists of 48% iron and lacks sulfur and phosphorus. Zinc, magnesium, and manganese commonly substitute for the iron, resulting in the siderite-smithsonite, siderite-magnesite, and siderite-rhodochrosite solid solution series.
Siderite has Mohs hardness of 3.75 to 4.25, a specific gravity of 3.96, a white streak and a vitreous lustre or pearly luster. Siderite is antiferromagnetic below its Néel temperature of which can assist in its identification.
It crystallizes in the trigonal crystal system, and are rhombohedral in shape, typically with curved and striated faces. It also occurs in masses. Color ranges from yellow to dark brown or black, the latter being due to the presence of manganese.
Siderite is commonly found in hydrothermal veins, and is associated with barite, fluorite, galena, and others. It is also a common diagenetic mineral in shales and sandstones, where it sometimes forms concretions, which can encase three-dimensionally preserved fossils. In sedimentary rocks, siderite commonly forms at shallow burial depths and its elemental composition is often related to the depositional environment of the enclosing sediments. In addition, a number of recent studies have used the oxygen isotopic composition of sphaerosiderite (a type associated with soils) as a proxy for the isotopic composition of meteoric water shortly after deposition.
Carbonate iron ore
Although carbonate iron ores, such as siderite, have been economically important for steel production, they are far from ideal as an ore.
Their hydrothermal mineralisation tends to form them as small ore lenses, often following steeply dipping bedding planes. This makes them not amenable to opencast working, and increases the cost of working them by mining with horizontal stopes. As the individual ore bodies are small, it may also be necessary to duplicate or relocate the pit head machinery, winding engine and pumping engine, between these bodies as each is worked out. This makes mining the ore an expensive proposition compared to typical ironstone or haematite opencasts.
The recovered ore also has drawbacks. The carbonate ore is more difficult to smelt than a haematite or other oxide ore. Driving off the carbonate as carbon dioxide requires more energy and so the ore 'kills' the blast furnace if added directly. Instead the ore must be given a preliminary roasting step. Developments of specific techniques to deal with these ores began in the early 19th century, largely with the work of Sir Thomas Lethbridge in Somerset. His 'Iron Mill' of 1838 used a three-chambered concentric roasting furnace, before passing the ore to a separate reducing furnace for smelting. Details of this mill were the invention of Charles Sanderson, a steel maker of Sheffield, who held the patent for it.
These differences between spathic ore and haematite have led to the failure of a number of mining concerns, notably the Brendon Hills Iron Ore Company.
Spathic iron ores are rich in manganese and have negligible phosphorus. This led to their one major benefit, connected with the Bessemer steel-making process. Although the first demonstrations by Bessemer in 1856 were successful, others' initial attempts to replicate his method infamously failed to produce good steel. Work by the metallurgist Robert Forester Mushet showed that the reason for the discrepancy was the nature of the Swedish ores that Bessemer had innocently used; they were very low in phosphorus. Using a typical European high-phosphorus ore in Bessemer's converter gave a poor quality steel. To produce high quality steel from a high-phosphorus ore, Mushet realised that he could operate the Bessemer converter for longer, burning off all the steel's impurities including the unwanted phosphorus but also the carbon (which is an essential ingredient in steel), and then re-adding carbon, along with manganese, in the form of a previously obscure ferromanganese ore with no phosphorus, spiegeleisen. This created a sudden demand for spiegeleisen. Although it was not available in sufficient quantity as a mineral, steelworks such as that at Ebbw Vale in South Wales soon learned to make it from the spathic siderite ores. For a few decades, spathic ores were therefore in demand and this encouraged their mining. In time though, the original 'acidic' liner of the Bessemer converter, made from siliceous sandstone or ganister, was replaced by a 'basic' liner in the newer Gilchrist Thomas process. This removed the phosphorus impurities as slag produced by chemical reaction with the liner, and no longer required spiegeleisen. From the 1880s demand for the ores fell once again and many of their mines, including those of the Brendon Hills, closed soon after.
Gallery
| Physical sciences | Minerals | Earth science |
875533 | https://en.wikipedia.org/wiki/Wollastonite | Wollastonite | Wollastonite is a calcium inosilicate mineral (CaSiO3) that may contain small amounts of iron, magnesium, and manganese substituting for calcium. It is usually white. It forms when impure limestone or dolomite is subjected to high temperature and pressure, which sometimes occurs in the presence of silica-bearing fluids as in skarns or in contact with metamorphic rocks. Associated minerals include garnets, vesuvianite, diopside, tremolite, epidote, plagioclase feldspar, pyroxene and calcite. It is named after the English chemist and mineralogist William Hyde Wollaston (1766–1828).
Despite its chemical similarity to the compositional spectrum of the pyroxene group of minerals—where magnesium (Mg) and iron (Fe) substitution for calcium ends with diopside and hedenbergite respectively—it is structurally very different, with a third tetrahedron in the linked chain (as opposed to two in the pyroxenes).
Production trends
Estimated world production of crude wollastonite ore was 1,200,000 tonnes in 2021. World reserves of wollastonite are estimated to exceed 100 million tonnes, though some existing deposits have not been surveyed.
Major producers of wollastonite include China, India, the United States, Mexico, and Finland.
In the United States, wollastonite is mined in Willsboro, New York (the first laboratory for local wollastonite research was in Essex, New York by Koert Burnham in the 1940s. The original laboratory building still exists as a residential & commercial building) and Gouverneur, New York. Deposits have also been mined commercially in North Western Mexico.
The price of raw wollastonite in 2008 varied between US$80 and US$500 per tonne depending on the country and size and shape of the powder particles.
Uses
Wollastonite is among the fastest reacting silicates, but may have high costs associated with carbon storage. Addition of wollastonite to soil stimulates organic carbon mineralization.
Ceramics
Wollastonite has industrial importance in ceramics manufacturing as an additive.
In ceramics, wollastonite decreases shrinkage and gas evolution during firing, increases green and fired strength, maintains brightness during firing, permits fast firing, and reduces crazing, cracking, and glaze defects.
Construction
Wollastonite can serve as a substitute for asbestos in floor tiles, friction products, insulating board and panels, paint, plastics, and roofing products. Similar to asbestos, wollastonite is resistant to chemical attack, stable at high temperatures, and improves flexural and tensile strength in composites. In some industries, wollastonite is used in different percentages of impurities, such as its use as a fabricator of mineral wool insulation, or as an ornamental building material.
Wollastonite is used in a cement announced in 2019 which "reduces the overall carbon footprint in precast concrete by 70%."
Wollastonite has been studied for carbon mineralization for storage of carbon dioxide (CO2) according to the following reaction:
Metallurgy
In metallurgical applications, wollastonite serves as a flux for welding, a source for calcium oxide, a slag conditioner, and to protect the surface of molten metal during the continuous casting of steel.
Paint
As an additive in paint, wollastonite improves the durability of the paint film, acts as a pH buffer, improves its resistance to weathering, reduces gloss, reduces pigment consumption, and acts as a flatting and suspending agent.
Plastic
In plastics, wollastonite improves tensile and flexural strength, reduces resin consumption, and improves thermal and dimensional stability at elevated temperatures. Surface treatments are used to improve the adhesion between the wollastonite and the polymers to which it is added.
Plastics and rubber applications were estimated to account for 25% to 35% of U.S. sales in 2009, followed by ceramics with 20% to 25%; paint, 10% to 15%; metallurgical applications, 10% to 15%; friction products, 10% to 15%; and miscellaneous, 10% to 15%. Ceramic applications probably account for 30% to 40% of wollastonite sales worldwide, followed by polymers (plastics and rubber) with 30% to 35% of sales, and paint with 10% to 15% of sales. The remaining sales were for construction, friction products, and metallurgical applications.
Substitutes
The acicular nature of many wollastonite products allows it to compete with other acicular materials, such as ceramic fiber, glass fiber, steel fiber, and several organic fibers, such as aramid, polyethylene, polypropylene, and polytetrafluoroethylene in products where improvements in dimensional stability, flexural modulus, and heat deflection are sought.
Wollastonite also competes with several nonfibrous minerals or rocks, such as kaolin, mica, and talc, which are added to plastics to increase flexural strength, and such minerals as barite, calcium carbonate, gypsum, and talc, which impart dimensional stability to plastics.
In ceramics, wollastonite competes with carbonates, feldspar, lime, and silica as a source of calcium and silicon. Its use in ceramics depends on the formulation of the ceramic body and the firing method.
Composition
In a pure CaSiO3, each component forms nearly half of the mineral by weight: 48.3% of CaO and 51.7% of SiO2. In some cases, small amounts of iron (Fe), and manganese (Mn), and lesser amounts of magnesium (Mg) substitute for calcium (Ca) in the mineral formula (e.g., rhodonite). Wollastonite can form a series of solid solutions in the system CaSiO3-FeSiO3, or hydrothermal synthesis of phases in the system MnSiO3-CaSiO3.
Geologic occurrence
Wollastonite usually occurs as a common constituent of a thermally metamorphosed impure limestone, it also could occur when the silicon is due to metamorphism in contact altered calcareous sediments, or to contamination in the invading igneous rock. In most of these occurrences it is the result of the following reaction between calcite and silica with the loss of carbon dioxide:
CaCO3 + SiO2 → CaSiO3 + CO2
Wollastonite may also be produced in a diffusion reaction in skarn, it develops when limestone within a sandstone is metamorphosed by a dike, which results in the formation of wollastonite in the sandstone as a result of outward migration of Ca.
Structure
Wollastonite crystallizes triclinically in space group P with the lattice constants a = 7.94 Å, b = 7.32 Å, c = 7.07 Å; α = 90,03°, β = 95,37°, γ = 103,43° and six formula units per unit cell. Wollastonite was once classed structurally among the pyroxene group, because both of these groups have a ratio of Si:O = 1:3. In 1931, Warren and Biscoe showed that the crystal structure of wollastonite differs from minerals of the pyroxene group, and they classified this mineral within a group known as the pyroxenoids. It has been shown that the pyroxenoid chains are more kinked than those of pyroxene group, and exhibit longer repeat distance. The structure of wollastonite contains infinite chains of [SiO4] tetrahedra sharing common vertices, running parallel to the b-axis. The chain motif in wollastonite repeats after three tetrahedra, whereas in pyroxenes only two are needed. The repeat distance in the wollastonite chains is 7.32 Å and equals the length of the crystallographic b-axis.
Molten CaSiO3 maintains a tetrahedral SiO4 local structure at temperatures up to 2000 °C. The nearest neighbor Ca-O coordination decreases from 6.0(2) in the room temperature glass to 5.0(2) in the 1700 °C liquid, coincident with an increasing number of longer Ca-O neighbors.
| Physical sciences | Silicate minerals | Earth science |
875717 | https://en.wikipedia.org/wiki/Merchant%20ship | Merchant ship | A merchant ship, merchant vessel, trading vessel, or merchantman is a watercraft that transports cargo or carries passengers for hire. This is in contrast to pleasure craft, which are used for personal recreation, and naval ships, which are used for military purposes.
They come in myriad sizes and shapes, from inflatable dive boats in Hawaii, to 5,000-passenger casino vessels on the Mississippi River, to tugboats plying New York Harbor, to oil tankers and container ships at major ports, to passenger-carrying submarines in the Caribbean.
Many merchant ships operate under a "flag of convenience" from a country other than the home of the vessel's owners, such as Liberia and Panama, which have more favorable maritime laws than other countries.
The Greek merchant marine is the largest in the world. Today, the Greek fleet accounts for some 16 per cent of the world's tonnage; this makes it currently the largest single international merchant fleet in the world, albeit not the largest in history.
During wars, merchant ships may be used as auxiliaries to the navies of their respective countries, and are called upon to deliver military personnel and materiel.
History
Definitions
The term "commercial vessel" is defined by the United States Coast Guard as any vessel (i.e. boat or ship) engaged in commercial trade or that carries passengers for hire.
In English, the term "Merchant Navy" without further clarification is used to refer to the British Merchant Navy; the United States merchant fleet is known as the United States Merchant Marine.
Name prefixes
Merchant ships' names have a prefix to indicate which kind of vessel they are:
CS = Cable Ship/Cable layer
LNG = Gas carrier transporting liquefied natural gas
LPG = Gas carrier transporting liquefied petroleum gas
MFV = Motor Fishing Vessel
MS = Motorship
MSV = Motor Stand-by Vessel
MT = Motor Tanker or Motor Tug Boat
MV = Motor/Merchant Vessel
MY = Motor Yacht
NS = Nuclear Ship
RMS = Royal Mail Ship
RRS = Royal Research Ship
RV = Research Vessel
SS = Steam Ship
SV = Sailing Vessel (although these can be sub coded as type of sailing vessel)
Merchant ship categories
The UNCTAD review of maritime transport categorizes ships as: oil tankers, bulk (and combination) carriers, general cargo ships, container ships, and "other ships", which includes "liquefied petroleum gas carriers, liquefied natural gas carriers, parcel (chemical) tankers, specialized tankers, reefers, offshore supply, tugs, dredgers, cruise, ferries, other non-cargo". General cargo ships include "multi-purpose and project vessels and Roll-on/roll-off cargo".
Cargo ship
A cargo ship or freighter is any sort of ship or vessel that carries cargo, goods, and materials from one port to another. Thousands of cargo carriers ply the world's seas and oceans each year; they handle the bulk of international trade. Cargo ships are usually specially designed for the task, often being equipped with cranes and other mechanisms to load and unload, and come in all sizes.
Bulk carrier
A bulk carrier is a ship used to transport bulk cargo items such as iron ore, bauxite, coal, cement, grain and similar cargo. Bulk carriers can be recognized by large box-like hatches on deck, designed to slide outboard or fold fore-and-aft to enable access for loading or discharging cargo. The dimensions of bulk carriers are often determined by the ports and sea routes that they need to serve, and by the maximum width of the Panama Canal. Most lakes are too small to accommodate bulk carriers, but a large fleet of lake freighters has been plying the Great Lakes and St. Lawrence Seaway of North America for over a century.
Container ship
A container ship is a cargo ship that carries its cargo in standardized containers, in a technique called containerization. These ships are a common means of commercial intermodal freight transport.
Tanker
A tanker is a ship designed to transport liquids in bulk. Tankers can range in size from several hundred tons, designed to serve small harbours and coastal settlements, to several hundred thousand tons, with these being designed for long-range haulage. A wide range of products are carried by tankers, including:
hydrocarbon products such as oil, LPG, and LNG
chemicals, such as ammonia, chlorine, and styrene monomer
fresh water
wine
Different products require different handling and transport, thus special types of tankers have been built, such as chemical tankers, oil tankers, and gas carriers.
Among oil tankers, supertankers were designed for carrying oil around the Horn of Africa from the Middle East; the FSO Knock Nevis being the largest vessel in the world, a ULCC supertanker formerly known as Jahre Viking (Seawise Giant). It has a deadweight of 565,000 metric tons and length of about . The use of such large ships is in fact very unprofitable, due to the inability to operate them at full cargo capacity; hence, the production of supertankers has currently ceased. Today's largest oil tankers in comparison by gross tonnage are TI Europe, TI Asia, TI Oceania, which are the largest sailing vessels today. But even with their deadweight of 441,585 metric tons, sailing as VLCC most of the time, they do not use more than 70% of their total capacity.
Apart from pipeline transport, tankers are the only method for transporting large quantities of oil, although such tankers have caused large environmental disasters when sinking close to coastal regions, causing oil spills. See , Erika, Exxon Valdez, Prestige and for examples of tankers that have been involved in oil spills.
Coastal trading vessel
Coastal trading vessels are smaller ships that carry any category of cargo along coastal, rather than trans-oceanic, routes. Coasters are shallow-hulled ships used for trade between locations on the same island or continent. Their shallow hulls allow them to sail over reefs and other submerged navigation hazards, whereas ships designed for blue-water trade usually have much deeper hulls for better seakeeping.
Passenger ship
A passenger ship is a ship whose primary function is to carry passengers. The category does not include cargo vessels which have accommodations for limited numbers of passengers, such as the formerly ubiquitous twelve-passenger freighters in which the transport of passengers is secondary to the carriage of freight. The type does however include many classes of ships which are designed to transport substantial numbers of passengers as well as freight. Indeed, until recently virtually all ocean liners were able to transport mail, package freight and express, and other cargo in addition to passenger luggage, and were equipped with cargo holds and derricks, kingposts, or other cargo-handling gear for that purpose. Modern cruiseferries have car decks for lorries as well as the passengers' cars. Only in more recent ocean liners and in virtually all cruise ships has this cargo capacity been removed. A ferry is a boat or ship carrying passengers and sometimes their vehicles. Ferries are also used to transport freight (in lorries and sometimes unpowered freight containers) and even railroad cars (in the case of a train ferry).
| Technology | Maritime transport | null |
875900 | https://en.wikipedia.org/wiki/Body%20of%20water | Body of water | A body of water or waterbody is any significant accumulation of water on the surface of Earth or another planet. The term most often refers to oceans, seas, and lakes, but it includes smaller pools of water such as ponds, wetlands, or more rarely, puddles. A body of water does not have to be still or contained; rivers, streams, canals, and other geographical features where water moves from one place to another are also considered bodies of water.
Most are naturally occurring geographical features, but some are artificial. There are types that can be either. For example, most reservoirs are created by engineering dams, but some natural lakes are used as reservoirs. Similarly, most harbors are naturally occurring bays, but some harbors have been created through construction.
Bodies of water that are navigable are known as waterways. Some bodies of water collect and move water, such as rivers and streams, and others primarily hold water, such as lakes and oceans.
Bodies of water are affected by gravity, which is what creates the tidal effects. Moreso, the impact of climate change on water is likely to intensify as observed through the rising sea levels, water acidification and flooding. This means that climate change has pressure on water bodies.
Climate change significantly affects bodies of water through rising temperatures, altered precipitation patterns, and sea-level rise. Warmer temperatures lead to the melting of glaciers and polar ice, contributing to rising sea levels and affecting coastal ecosystems. Freshwater bodies, such as rivers and lakes, are experiencing more frequent droughts, affecting water availability for communities and biodiversity. Moreover, ocean acidification, caused by increased carbon dioxide absorption, threatens marine ecosystems like coral reefs. Collaborative global efforts are needed to mitigate these impacts through sustainable water management practices.
Types
Bodies of water can be categorized into:
Rain water
Surface water
Underground water
There are some geographical features involving water that are not bodies of water, for example, waterfalls, geysers and rapids.
Gallery
| Physical sciences | Hydrology | null |
876428 | https://en.wikipedia.org/wiki/Divergent%20series | Divergent series | In mathematics, a divergent series is an infinite series that is not convergent, meaning that the infinite sequence of the partial sums of the series does not have a finite limit.
If a series converges, the individual terms of the series must approach zero. Thus any series in which the individual terms do not approach zero diverges. However, convergence is a stronger condition: not all series whose terms approach zero converge. A counterexample is the harmonic series
The divergence of the harmonic series was proven by the medieval mathematician Nicole Oresme.
In specialized mathematical contexts, values can be objectively assigned to certain series whose sequences of partial sums diverge, in order to make meaning of the divergence of the series. A summability method or summation method is a partial function from the set of series to values. For example, Cesàro summation assigns Grandi's divergent series
the value . Cesàro summation is an averaging method, in that it relies on the arithmetic mean of the sequence of partial sums. Other methods involve analytic continuations of related series. In physics, there are a wide variety of summability methods; these are discussed in greater detail in the article on regularization.
History
Before the 19th century, divergent series were widely used by Leonhard Euler and others, but often led to confusing and contradictory results. A major problem was Euler's idea that any divergent series should have a natural sum, without first defining what is meant by the sum of a divergent series. Augustin-Louis Cauchy eventually gave a rigorous definition of the sum of a (convergent) series, and for some time after this, divergent series were mostly excluded from mathematics. They reappeared in 1886 with Henri Poincaré's work on asymptotic series. In 1890, Ernesto Cesàro realized that one could give a rigorous definition of the sum of some divergent series, and defined Cesàro summation. (This was not the first use of Cesàro summation, which was used implicitly by Ferdinand Georg Frobenius in 1880; Cesàro's key contribution was not the discovery of this method, but his idea that one should give an explicit definition of the sum of a divergent series.) In the years after Cesàro's paper, several other mathematicians gave other definitions of the sum of a divergent series, although these are not always compatible: different definitions can give different answers for the sum of the same divergent series; so, when talking about the sum of a divergent series, it is necessary to specify which summation method one is using.
Examples
1 - 1 + 1 - 1 + ⋯
1 − 2 + 3 − 4 + ⋯
1 − 1 + 2 − 6 + 24 − 120 + ⋯
1 − 2 + 4 − 8 + ⋯
1 + 2 + 4 + 8 + ⋯
1 + 1 + 1 + 1 + ⋯
1 + 2 + 3 + 4 + ⋯
Theorems on methods for summing divergent series
A summability method M is regular if it agrees with the actual limit on all convergent series. Such a result is called an Abelian theorem for M, from the prototypical Abel's theorem. More subtle, are partial converse results, called Tauberian theorems, from a prototype proved by Alfred Tauber. Here partial converse means that if M sums the series Σ, and some side-condition holds, then Σ was convergent in the first place; without any side-condition such a result would say that M only summed convergent series (making it useless as a summation method for divergent series).
The function giving the sum of a convergent series is linear, and it follows from the Hahn–Banach theorem that it may be extended to a summation method summing any series with bounded partial sums. This is called the Banach limit. This fact is not very useful in practice, since there are many such extensions, inconsistent with each other, and also since proving such operators exist requires invoking the axiom of choice or its equivalents, such as Zorn's lemma. They are therefore nonconstructive.
The subject of divergent series, as a domain of mathematical analysis, is primarily concerned with explicit and natural techniques such as Abel summation, Cesàro summation and Borel summation, and their relationships. The advent of Wiener's tauberian theorem marked an epoch in the subject, introducing unexpected connections to Banach algebra methods in Fourier analysis.
Summation of divergent series is also related to extrapolation methods and sequence transformations as numerical techniques. Examples of such techniques are Padé approximants, Levin-type sequence transformations, and order-dependent mappings related to renormalization techniques for large-order perturbation theory in quantum mechanics.
Properties of summation methods
Summation methods usually concentrate on the sequence of partial sums of the series. While this sequence does not converge, we may often find that when we take an average of larger and larger numbers of initial terms of the sequence, the average converges, and we can use this average instead of a limit to evaluate the sum of the series. A summation method can be seen as a function from a set of sequences of partial sums to values. If A is any summation method assigning values to a set of sequences, we may mechanically translate this to a series-summation method AΣ that assigns the same values to the corresponding series. There are certain properties it is desirable for these methods to possess if they are to arrive at values corresponding to limits and sums, respectively.
Regularity. A summation method is regular if, whenever the sequence s converges to x, Equivalently, the corresponding series-summation method evaluates
Linearity. A is linear if it is a linear functional on the sequences where it is defined, so that for sequences r, s and a real or complex scalar k. Since the terms of the series a are linear functionals on the sequence s and vice versa, this is equivalent to AΣ being a linear functional on the terms of the series.
Stability (also called translativity). If s is a sequence starting from s0 and s′ is the sequence obtained by omitting the first value and subtracting it from the rest, so that , then A(s) is defined if and only if A(s′) is defined, and Equivalently, whenever for all n, then Another way of stating this is that the shift rule must be valid for the series that are summable by this method.
The third condition is less important, and some significant methods, such as Borel summation, do not possess it.
One can also give a weaker alternative to the last condition.
Finite re-indexability. If a and a′ are two series such that there exists a bijection such that for all i, and if there exists some such that for all i > N, then (In other words, a′ is the same series as a, with only finitely many terms re-indexed.) This is a weaker condition than stability, because any summation method that exhibits stability also exhibits finite re-indexability, but the converse is not true.)
A desirable property for two distinct summation methods A and B to share is consistency: A and B are consistent if for every sequence s to which both assign a value, (Using this language, a summation method A is regular iff it is consistent with the standard sum Σ.) If two methods are consistent, and one sums more series than the other, the one summing more series is stronger.
There are powerful numerical summation methods that are neither regular nor linear, for instance nonlinear sequence transformations like Levin-type sequence transformations and Padé approximants, as well as the order-dependent mappings of perturbative series based on renormalization techniques.
Taking regularity, linearity and stability as axioms, it is possible to sum many divergent series by elementary algebraic manipulations. This partly explains why many different summation methods give the same answer for certain series.
For instance, whenever the geometric series
can be evaluated regardless of convergence. More rigorously, any summation method that possesses these properties and which assigns a finite value to the geometric series must assign this value. However, when r is a real number larger than 1, the partial sums increase without bound, and averaging methods assign a limit of infinity.
Classical summation methods
The two classical summation methods for series, ordinary convergence and absolute convergence, define the sum as a limit of certain partial sums. These are included only for completeness; strictly speaking they are not true summation methods for divergent series since, by definition, a series is divergent only if these methods do not work. Most but not all summation methods for divergent series extend these methods to a larger class of sequences.
Absolute convergence
Absolute convergence defines the sum of a sequence (or set) of numbers to be the limit of the net of all partial sums , if it exists. It does not depend on the order of the elements of the sequence, and a classical theorem says that a sequence is absolutely convergent if and only if the sequence of absolute values is convergent in the standard sense.
Sum of a series
Cauchy's classical definition of the sum of a series defines the sum to be the limit of the sequence of partial sums . This is the default definition of convergence of a sequence.
Nørlund means
Suppose pn is a sequence of positive terms, starting from p0. Suppose also that
If now we transform a sequence s by using p to give weighted means, setting
then the limit of tn as n goes to infinity is an average called the Nørlund mean Np(s).
The Nørlund mean is regular, linear, and stable. Moreover, any two Nørlund means are consistent.
Cesàro summation
The most significant of the Nørlund means are the Cesàro sums. Here, if we define the sequence pk by
then the Cesàro sum Ck is defined by Cesàro sums are Nørlund means if , and hence are regular, linear, stable, and consistent. C0 is ordinary summation, and C1 is ordinary Cesàro summation. Cesàro sums have the property that if then Ch is stronger than Ck.
Abelian means
Suppose } is a strictly increasing sequence tending towards infinity, and that . Suppose
converges for all real numbers x > 0. Then the Abelian mean Aλ is defined as
More generally, if the series for f only converges for large x but can be analytically continued to all positive real x, then one can still define the sum of the divergent series by the limit above.
A series of this type is known as a generalized Dirichlet series; in applications to physics, this is known as the method of heat-kernel regularization.
Abelian means are regular and linear, but not stable and not always consistent between different choices of λ. However, some special cases are very important summation methods.
Abel summation
If , then we obtain the method of Abel summation. Here
where z = exp(−x). Then the limit of f(x) as x approaches 0 through positive reals is the limit of the power series for f(z) as z approaches 1 from below through positive reals, and the Abel sum A(s) is defined as
Abel summation is interesting in part because it is consistent with but more powerful than Cesàro summation: whenever the latter is defined. The Abel sum is therefore regular, linear, stable, and consistent with Cesàro summation.
Lindelöf summation
If , then (indexing from one) we have
Then L(s), the Lindelöf sum, is the limit of f(x) as x goes to positive zero. The Lindelöf sum is a powerful method when applied to power series among other applications, summing power series in the Mittag-Leffler star.
If g(z) is analytic in a disk around zero, and hence has a Maclaurin series G(z) with a positive radius of convergence, then in the Mittag-Leffler star. Moreover, convergence to g(z) is uniform on compact subsets of the star.
Analytic continuation
Several summation methods involve taking the value of an analytic continuation of a function.
Analytic continuation of power series
If Σanxn converges for small complex x and can be analytically continued along some path from x = 0 to the point x = 1, then the sum of the series can be defined to be the value at x = 1. This value may depend on the choice of path. One of the first examples of potentially different sums for a divergent series, using analytic continuation, was given by Callet, who observed that if then
Evaluating at , one gets
However, the gaps in the series are key. For for example, we actually would get
, so different sums correspond to different placements of the 's.
Another example of analytic continuation is the divergent alternating series
which is a sum over products of -functions and Pochhammer's symbols. Using the
duplication formula of the -function, it reduces to
a generalized hypergeometric series
Euler summation
Euler summation is essentially an explicit form of analytic continuation. If a power series converges for small complex z and can be analytically continued to the open disk with diameter from to 1 and is continuous at 1, then its value at q is called the Euler or (E,q) sum of the series Σan. Euler used it before analytic continuation was defined in general, and gave explicit formulas for the power series of the analytic continuation.
The operation of Euler summation can be repeated several times, and this is essentially equivalent to taking an analytic continuation of a power series to the point z = 1.
Analytic continuation of Dirichlet series
This method defines the sum of a series to be the value of the analytic continuation of the Dirichlet series
at s = 0, if this exists and is unique. This method is sometimes confused with zeta function regularization.
If s = 0 is an isolated singularity, the sum is defined by the constant term of the Laurent series expansion.
Zeta function regularization
If the series
(for positive values of the an) converges for large real s and can be analytically continued along the real line to s = −1, then its value at s = −1 is called the zeta regularized sum of the series a1 + a2 + ... Zeta function regularization is nonlinear. In applications, the numbers ai are sometimes the eigenvalues of a self-adjoint operator A with compact resolvent, and f(s) is then the trace of A−s. For example, if A has eigenvalues 1, 2, 3, ... then f(s) is the Riemann zeta function, ζ(s), whose value at s = −1 is −, assigning a value to the divergent series . Other values of s can also be used to assign values for the divergent sums , and in general
where Bk is a Bernoulli number.
Integral function means
If J(x) = Σpnxn is an integral function, then the J sum of the series a0 + ... is defined to be
if this limit exists.
There is a variation of this method where the series for J has a finite radius of convergence r and diverges at x = r. In this case one defines the sum as above, except taking the limit as x tends to r rather than infinity.
Borel summation
In the special case when J(x) = ex this gives one (weak) form of Borel summation.
Valiron's method
Valiron's method is a generalization of Borel summation to certain more general integral functions J. Valiron showed that under certain conditions it is equivalent to defining the sum of a series as
where H is the second derivative of G and c(n) = e−G(n), and a0 + ... + ah is to be interpreted as 0 when h < 0.
Moment methods
Suppose that dμ is a measure on the real line such that all the moments
are finite. If a0 + a1 + ... is a series such that
converges for all x in the support of μ, then the (dμ) sum of the series is defined to be the value of the integral
if it is defined. (If the numbers μn increase too rapidly then they do not uniquely determine the measure μ.)
Borel summation
For example, if dμ = e−x dx for positive x and 0 for negative x then μn = n!, and this gives one version of Borel summation, where the value of a sum is given by
There is a generalization of this depending on a variable α, called the (B′,α) sum, where the sum of a series a0 + ... is defined to be
if this integral exists. A further generalization is to replace the sum under the integral by its analytic continuation from small t.
Miscellaneous methods
BGN hyperreal summation
This summation method works by using an extension to the real numbers known as the hyperreal numbers. Since the hyperreal numbers include distinct infinite values, these numbers can be used to represent the values of divergent series. The key method is to designate a particular infinite value that is being summed, usually , which is used as a unit of infinity. Instead of summing to an arbitrary infinity (as is typically done with ), the BGN method sums to the specific hyperreal infinite value labeled . Therefore, the summations are of the form
This allows the usage of standard formulas for finite series such as arithmetic progressions in an infinite context. For instance, using this method, the sum of the progression is , or, using just the most significant infinite hyperreal part, .
Hausdorff transformations
.
Hölder summation
Hutton's method
In 1812 Hutton introduced a method of summing divergent series by starting with the sequence of partial sums, and repeatedly applying the operation of replacing a sequence s0, s1, ... by the sequence of averages , , ..., and then taking the limit.
Ingham summability
The series a1 + ... is called Ingham summable to s if
Albert Ingham showed that if δ is any positive number then (C,−δ) (Cesàro) summability implies Ingham summability, and Ingham summability implies (C,δ) summability.
Lambert summability
The series a1 + ... is called Lambert summable to s if
If a series is (C,k) (Cesàro) summable for any k then it is Lambert summable to the same value, and if a series is Lambert summable then it is Abel summable to the same value.
Le Roy summation
The series a0 + ... is called Le Roy summable to s if
Mittag-Leffler summation
The series a0 + ... is called Mittag-Leffler (M) summable to s if
Ramanujan summation
Ramanujan summation is a method of assigning a value to divergent series used by Ramanujan and based on the Euler–Maclaurin summation formula. The Ramanujan sum of a series f(0) + f(1) + ... depends not only on the values of f at integers, but also on values of the function f at non-integral points, so it is not really a summation method in the sense of this article.
Riemann summability
The series a1 + ... is called (R,k) (or Riemann) summable to s if
The series a1 + ... is called R2 summable to s if
Riesz means
If λn form an increasing sequence of real numbers and
then the Riesz (R,λ,κ) sum of the series a0 + ... is defined to be
Vallée-Poussin summability
The series a1 + ... is called VP (or Vallée-Poussin) summable to s if
where is the gamma function.
Zeldovich summability
The series is Zeldovich summable if
| Mathematics | Sequences and series | null |
876719 | https://en.wikipedia.org/wiki/DPT%20vaccine | DPT vaccine | The DPT vaccine or DTP vaccine is a class of combination vaccines to protect
against three infectious diseases in humans: diphtheria, pertussis (whooping cough), and tetanus (lockjaw). The vaccine components include diphtheria and tetanus toxoids, and either killed whole cells of the bacterium that causes pertussis or pertussis antigens. The term toxoid refers to vaccines which use an inactivated toxin produced by the pathogen which they are targeted against to generate an immune response. In this way, the toxoid vaccine generates an immune response which is targeted against the toxin which is produced by the pathogen and causes disease, rather than a vaccine which is targeted against the pathogen itself. The whole cells or antigens will be depicted as either "DTwP" or "DTaP", where the lower-case "w" indicates whole-cell inactivated pertussis and the lower-case "a" stands for "acellular". In comparison to alternative vaccine types, such as live attenuated vaccines, the DTP vaccine does not contain any live pathogen, but rather uses inactivated toxoid (and for pertussis, either a dead pathogen or pure antigens) to generate an immune response; therefore, there is not a risk of use in populations that are immune compromised since there is not any known risk of causing the disease itself. As a result, the DTP vaccine is considered a safe vaccine to use in anyone and it generates a much more targeted immune response specific for the pathogen of interest.
In the United States, the DPT (whole-cell) vaccine was administered as part of the childhood vaccines recommended by the Centers for Disease Control and Prevention (CDC) until 1996, when the acellular DTaP vaccine was licensed for use.
History
Diphtheria and tetanus toxoids and whole-cell pertussis (DTP; now also "DTwP" to differentiate from the broader class of triple-combination vaccines) vaccination was licensed in 1949. Since the introduction of the combination vaccine, there has been an extensive decline in the incidence of pertussis, or whooping cough, the disease which the vaccine protects against. Additionally, the rates of disease have continued to decline as more extensive immunization strategies have been implemented, including booster doses and increased emphasis on increasing health literacy.
In the 20th century, the advancements in vaccinations helped to reduce the incidence of childhood pertussis and had a dramatically positive effect on the health of populations in the United States. However, in the early 21st century, reported instances of the disease increased 20-fold due to a downturn in the number of immunizations received and resulted in numerous fatalities. During the 21st century, many parents declined to vaccinate their children against pertussis for fear of perceived side effects despite scientific evidence showing vaccines to be highly effective and safe. In 2009, the journal Pediatrics concluded the largest risk among unvaccinated children was not the contraction of side effects, but rather the disease that the vaccination aims to protect against.
DTP vaccines with acellular pertussis (DTaP; see below) were introduced in the 1990s. The reduced range of antigens causes fewer side effects, but results in a more expensive, shorter-lasting, and possibly less protective vaccine compared to DTwP. High-income countries have mostly switched to DTaP. As of 2023, global production of aP remains limited.
Vaccination rates
In 2016, the CDC reported that 80.4% of children in the US have received four or more DTaP vaccinations by 2 years of life. Vaccination rates for children aged 13–17 with one or more TDaP shots was 90.2% in 2019. Only 43.6% of adults (older than 18) have received a TDaP shot in the last 10 years. The CDC aims to increase vaccination rate among 2-year-olds from 80.4% to 90.0%
The World Health Organization (WHO) estimates that 89% of people globally have received at least one dose of DTP vaccine and 84% have received three doses of the vaccine, completing the WHO-recommended primary series (DTP3). The WHO also tracks the DTP3 completion rate among one-year-olds on a yearly basis. Yearly DTP3 completion rate is considered a good proxy for the completeness of childhood vaccination in general.
Combination vaccines with acellular pertussis
DTaP and Tdap are both combination vaccines against diphtheria, tetanus, and pertussis. The "a" indicates that the pertussis toxoids are acellular, while the lower-case "d" and "p" in "Tdap" indicate smaller concentrations of diphtheria toxoids and pertussis antigens.
DTaP
DTaP (also DTP and TDaP) is a combination vaccine against diphtheria, tetanus, and pertussis, in which the pertussis component is acellular. This is in contrast to whole-cell, inactivated DTP (or DTwP). The acellular vaccine uses selected antigens of the pertussis pathogen to induce immunity. Because it uses fewer antigens than the whole-cell vaccines, it is considered to cause fewer side effects, but it is also more expensive. Research suggests that the DTwP vaccine is more effective than DTaP in conferring immunity, because DTaP's narrower antigen base is less effective against current pathogen strains.
Tdap
Tdap (also TDP) is a tetanus toxoid, reduced diphtheria toxoid, and acellular pertussis vaccine. It was licensed in the United States for use in adults and adolescents on 10 June 2005. Two Tdap vaccines are available in the US. In January 2011, the US Centers for Disease Control and Prevention (CDC) Advisory Committee on Immunization Practices (ACIP) recommended the use of Tdap in adults of all ages, including those age 65 and above. In October 2011, in an effort to reduce the burden of pertussis in infants, the ACIP recommended that unvaccinated pregnant women receive a dose of Tdap. On 24 October 2012, the ACIP voted to recommend the use of Tdap during every pregnancy.
The ACIP and Canada's National Advisory Committee on Immunization (NACI) recommended that both adolescents and adults receive Tdap in place of their next Td booster (recommended to be given every ten years). Tdap and Td can be used as prophylaxis for tetanus in wound management. People who will be in contact with young infants are encouraged to get Tdap even if it has been less than five years since Td or TT to reduce the risk of infants being exposed to pertussis. NACI suggests intervals shorter than five years can be used for catch-up programs and other instances where programmatic concerns make five-year intervals difficult.
The WHO recommends a pentavalent vaccine, combining the DTP vaccine with vaccines against Haemophilus influenzae type B and hepatitis B. Evidence on how effective this pentavalent vaccine is compared to the individual vaccines has not yet been determined.
A 2019 study found that state requirements mandating the use of the Tdap vaccine "increased Tdap vaccine take-up and reduced pertussis (whooping cough) incidence by about 32%."
Related combination vaccines
Excluding pertussis
DT and Td vaccines lack the pertussis component. The Td vaccine is administered to children over the age of seven as well as to adults. It is most commonly administered as a booster shot every 10 years. The Td booster shot may also be administered as protection from a severe burn or dirty wound. The DT vaccine is given to children under the age of seven who are unable to receive the pertussis antigen in the DTaP vaccine due to a contraindication.
Additional targets
In the United States, a combined inactivated polio (IPV), DTaP, and hepatitis B DTaP-IPV-HepB vaccine is available for children. In the UK, all babies born on or after 1 August 2017 are offered a hexavalent vaccine: DTaP, IPV, Haemophilus influenzae, and hepatitis B (DTaP-Hib-HepB-IPV in short).
As of 2023, most of the DTP vaccine procured by UNICEF is of the DTwP-HepB-Hib (pentavalent whole-cell) type. The UNICEF plans to procure the DTwP-HepB-Hib-IPV (hexavalent whole-cell) vaccine starting in 2024.
Contraindications
The DPT vaccine should be avoided in persons who experienced a severe allergic reaction, such as anaphylaxis, to a past vaccine containing tetanus, diphtheria, or pertussis. It should also be avoided in persons with a known severe allergy to an ingredient in the vaccine. If the reaction was caused by tetanus toxoids, the CDC recommends considering a passive immunization with tetanus immune globulin (TIG) if a person has a large or unclean wound. The DPT vaccine should also be avoided if a person developed encephalopathy (seizures, coma, declined consciousness) within seven days of receiving any pertussis-containing vaccine and the encephalopathy cannot be traced to another cause. A DT vaccine is available for children under the ages of seven who have contraindications or precautions to pertussis-containing vaccines.
Side effects
DTaP
Common side effects include soreness where the shot was given, fever, irritability, tenderness, loss of appetite, and vomiting. Most side effects are mild to moderate and may last from one to three days. More serious but rare reactions after a DTaP vaccination may include seizures, lowered consciousness, or a high fever over . Allergic reactions are uncommon, but are medical emergencies. Signs of an allergic reaction include hives, dyspnea, wheezing, swelling of face and throat, syncope, and tachycardia and the child should be rushed to the nearest hospital.
Tdap
Common side effects include pain or swelling where the shot was given, mild fever, headache, tiredness, nausea, vomiting, diarrhea, and stomach ache. Allergic reactions are possible and have the same presentation and indications as described above for allergic reactions in DTaP. Any individual who has experienced a life-threatening allergic reaction after receiving a previous dose of diphtheria, tetanus, or pertussis containing vaccine should not receive the Tdap vaccination.
In pregnant women, research suggests that Tdap administration may be associated with an increased risk of chorioamnionitis, a placental infection. Increased incidence of fever is also noted in pregnant women. Despite the observed increase in incidence of chorioamnionitis in pregnant women following Tdap administration, there has been no observed increase in the incidence of preterm birth, for which chorioamnionitis is a risk factor. Research has not discerned an association between Tdap administration during pregnancy and other serious pregnancy complications such as neonatal death and stillbirth. An association between Tdap administration during pregnancy and pregnancy-related hypertensive disorders (such as pre-eclampsia) has not been identified.
Immunization schedules and requirements
France
In France, children are given DTaP-Hib-HepB-IPV vaccines at 2 months (first dose) and 4 months (second dose) with a booster at 11 months of age. A tetravalent booster for diphtheria, pertussis, tetanus and poliomyelitis is given at 6 years, at 11–13 years, then at 25, 45, 65 years of age, then every 10 years.
Netherlands
In the Netherlands, pertussis is known as kinkhoest and DKTP refers to the DTaP-IPV combination vaccine against diphtheria, kinkhoest, tetanus, and polio. DTaP is given as part of the National Immunisation Programme.
United Kingdom
In the United Kingdom, Td/IPV is called the "3-in-1 teenage booster" and protects against tetanus, diphtheria and polio. It is given by the NHS to all teenagers aged 14 (the hexavalent vaccine is given to infants and provides the first stage of protection against diphtheria, tetanus, and polio, as well as pertussis, Haemophilus influenzae type B and hepatitis B). Subsequent boosters are recommended for foreign travellers where more than 10 years has passed since their last booster. This is provided on the NHS free of charge due to the significant risk that an imported case of polio could pose to public health in Britain.
United States
The standard immunization regimen for children within the United States is five doses of DTaP between the ages of two months and fifteen years. To be considered fully vaccinated, the Centers for Disease Control and Prevention (CDC) typically requires five doses of Tdap. The CDC recommends that children receive their first dose at two months, the second dose at four months, the third dose at six months, the fourth dose between 15 and 18 months, and the fifth dose between 4–6 years. If the fourth dose of the DTaP immunization regimen falls on or subsequent to the recipient's fourth birthday, the CDC states that only four doses are required to be fully vaccinated. In the instance that an individual under 18 has not received the DTaP vaccine, individuals should be vaccinated on the schedule in accordance with the vaccination "catch up schedule" provided by the CDC.
Infants younger than twelve months of age, specifically less than three months of age, are at highest risk of acquiring pertussis. In U.S., there is no current tetanus-diphtheria-pertussis vaccination (whooping cough) recommended or licensed for new born infants. As a result, in their first few months of life, unprotected infants are at highest risk of life-threatening complications and infections from pertussis. Infants should not receive pertussis vaccination younger than six weeks of age. Ideally, Infants should receive DTaP (name of whooping cough vaccine for children from age 2 months through 6 years) at 2, 4, 6 months of age and they are not protected until the full series is completed. To protect infants younger than twelve months of age not vaccinated with Tdap against pertussis, ACIP also recommends adults (e.g., parents, siblings, grandparents, childcare providers, and healthcare personnel) and children to receive Tdap at least two weeks before being in contact with the infant.
The CDC recommends that adults who have received their childhood DTP series receive a Td or Tdap booster every ten years. For adults that have not received the DTP series, the CDC recommends a three-part vaccine series followed by a Td or Tdap booster every ten years.
In pregnancy
According to the CDC's Advisory Committee on Immunization Practices (ACIP) guidelines, one dose of Tdap is recommended during each pregnancy to ensure protection against pertussis in newborn infants. Optimal timing to administer a dose of Tdap during each pregnancy is between 27 through 36 weeks gestation. If Tdap is administered early in pregnancy, it is not recommended to administer again during the 27 through 36 weeks gestation period as only one dose is recommended during pregnancy. In October 2022, Boostrix (Tetanus Toxoid, Reduced Diphtheria Toxoid and Acellular Pertussis Vaccine, Adsorbed [Tdap]) was approved for immunization during the third trimester of pregnancy to prevent pertussis, commonly known as whooping cough, in infants younger than two months of age.
Pregnant women who have not previously vaccinated with Tdap (i.e., have never received DTP, DTaP, or DT as child or Td or TT as an adult) are recommended to receive a series of three Td vaccinations starting during pregnancy to ensure protection against maternal and neonatal tetanus. In such cases, administration of Tdap is recommended after 20 weeks' gestation, and in earlier pregnancy a single dose of Tdap can be substituted for one dose of Td, and then the series completed with Td. For pregnant women not previously vaccinated with Tdap, if Tdap is not administered during pregnancy, it should be administered immediately postpartum. Postpartum administration of TDaP is not equivalent to administration of the vaccination during pregnancy. Because the vaccine is administered postpartum, the mother is unable to develop antibodies that can be transferred to the infant in utero, consequently, leaving the infant vulnerable to the diseases preventable by the Tdap Vaccine. Postpartum administration of the TdaP vaccine to the mother seeks to reduce the likelihood that the mother will contract disease that can be subsequently passed on the infant, albeit there will still be a two-week period prior to the protective effects of the vaccine setting in. Postpartum administration is an extension of the concept of "cocooning", a term that refers to the full vaccination of all individuals that may come into direct contact with the infant. Cocooning, like postpartum Tdap administration, is not recommended by the CDC. Cocooning depends on ensuring full vaccination of all individuals that the infant may come into contact with, and there may be financial, administrative or personal barriers that preclude full and timely vaccination of all individuals within the "cocoon".
Brand names
Australia
United Kingdom
Brand names in the United Kingdom include Revaxis (Sanofi Pasteur).
United States
, there are six DTaP vaccines and two Tdap vaccines licensed and available for use in the United States. All of them are indicated as childhood vaccinations with the schedules as follows:
| Biology and health sciences | Vaccines | Health |
877127 | https://en.wikipedia.org/wiki/Fastener | Fastener | A fastener (US English) or fastening (UK English) is a hardware device that mechanically joins or affixes two or more objects together. In general, fasteners are used to create non-permanent joints; that is, joints that can be removed or dismantled without damaging the joining components. Steel fasteners are usually made of stainless steel, carbon steel, or alloy steel.
Other methods of joining materials, some of which may create permanent joints, include: crimping, welding, soldering, brazing, taping, gluing, cement, or the use of other adhesives. Force may also be used, such as with magnets, vacuum (like suction cups), or even friction (like sticky pads). Some types of woodworking joints make use of separate internal reinforcements, such as dowels or biscuits, which in a sense can be considered fasteners within the scope of the joint system, although on their own they are not general-purpose fasteners.
Furniture supplied in flat-pack form often uses cam dowels locked by cam locks, also known as conformat fasteners. Fasteners can also be used to close a container such as a bag, a box, or an envelope; or they may involve keeping together the sides of an opening of flexible material, attaching a lid to a container, etc. There are also special-purpose closing devices, e.g., a bread clip.
Items like a rope, string, wire, cable, chain, or plastic wrap may be used to mechanically join objects; however, because they have additional common uses, they are not generally categorized as fasteners. Likewise, hinges and springs may join objects together, but they are ordinarily not considered fasteners because their primary purpose is to allow articulation rather than rigid affixment.
Industry
In 2005, it was estimated that the United States fastener industry runs 350 manufacturing plants and employs 40,000 workers. The industry is strongly tied to the production of automobiles, aircraft, appliances, agricultural machinery, commercial construction, and infrastructure. More than 200 billion fasteners are used per year in the U.S., 26 billion of these by the automotive industry. The largest distributor of fasteners in North America is the Fastenal Company.
Materials
There are three major steel fasteners used in industries: stainless steel, carbon steel, and alloy steel. The major grade used in stainless steel fasteners: 200 series, 300 series, and 400 series. Titanium, aluminium, and various alloys are also common materials of construction for metal fasteners. In many cases, special coatings or plating may be applied to metal fasteners to improve their performance characteristics by, for example, enhancing corrosion resistance. Common coatings/platings include zinc, chrome, and hot-dip galvanizing.
Applications
When selecting a fastener for industrial applications, it is important to consider a variety of factors. The threading, the applied load on the fastener, the stiffness of the fastener, and the number of fasteners needed should all be taken into account.
When choosing a fastener for a given application, it is important to know the specifics of that application to help select the proper material for the intended use. Factors that should be considered include:
Accessibility
Environment, including temperature, water exposure, and potentially corrosive elements
Installation process
Materials to be joined
Reusability
Weight restrictions
Types
A threaded fastener has internal or external screw threads. The most common types are the screw, nut and bolt, possibly involving washers.
Other more specialized types of threaded fasteners include captive threaded fasteners, stud, threaded inserts, and threaded rods.
Other types of fastener include:
anchor bolt
batten
bolt (fastener)
screw
bolt snap
brass fastener
buckle
button
cable tie
cam
captive fastener
clamp (or cramp)
hose clamp
clasp and shackle
bolt snap
carabiner
circle cotter
lobster clasp
cleco
clip
Binder clip
Bulldog clip
Crocodile clip
circlip
Clothespin
hairpin clip
paper clip
terry clip
clutch
drawing pin (thumbtack)
flange
frog
grommet
hook-and-eye closure
hook and loop fastener
Velcro
latch
nail and rivet
solid/round head rivets
semi-tubular rivets
blind (pop) rivet
pegs
tent peg
PEM nut
pins
clevis fastener
cotter
dowel
linchpin
R-clip
safety pin
split pin
spring pin
tapered pin
retaining rings
circlip
e-ring
rivet-like
well nut
rock bolt
rubber band (or bands of other materials)
screw anchor
snap fastener
snap-fit
staple
stitches
strap
tie
toggle bolt
tolerance rings
treasury tag
twist tie
wedge anchor
zipper
Common fastener head styles
Common head styles include:
Flat head fasteners: Ideal for applications where aesthetics are a priority, flat head fasteners sit flush with the surface, offering a clean appearance.
Round head fasteners: With a rounded top, round head fasteners provide a larger bearing surface, suitable for sheet metal or thin plastic assemblies.
Pan head fasteners: Pan head fasteners combine a slightly flattened top with a larger bearing surface, offering a streamlined appearance for aesthetic applications.
Socket head fasteners: Designed for high torque applications, socket head fasteners are driven with a hex key, reducing the risk of cam-out.
Hex head fasteners: Known for their high torque capacity, hex head fasteners are easily driven with a spanner or wrench, ideal for heavy-duty applications.
Square head fasteners: Offering increased wrenching area and reduced risk of rounding off, square head fasteners are used in high torque applications.
Flange head fasteners: Integrating a flange for a larger bearing surface, flange head fasteners distribute clamping force without damaging the material.
Wing head fasteners: Featuring protruding "wings" for hand tightening, wing head fasteners are suitable for applications requiring frequent adjustments.
T-slot fasteners: Designed for T-slotted aluminium extrusions, T-slot fasteners provide a secure and adjustable connection for framing and guarding systems.
Standards & traceability
There are multiple standards bodies for fasteners, including the US Industrial Fasteners Institute and the European Industrial Fastener Institute.
ASME B18 standards on certain fasteners
The American Society of Mechanical Engineers (ASME) publishes several standards on fasteners. Some are:
B18.3 Socket Cap, Shoulder, Set Screws, and Hex Keys (Inch Series)
B18.6.1 Wood Screws (Inch Series)
B18.6.2 Slotted Head Cap Screws, Square Head Set Screws, And Slotted Headless Set Screws (Inch Series)
B18.6.3 Machine Screws, Tapping Screws, and Metallic Drive Screws (Inch Series)
B18.18 Quality Assurance For Fasteners
B18.24 Part Identifying Number (PIN) Code System Standard for B18 Fastener Products
For military hardware
American screws, bolts, and nuts were historically not fully interchangeable with their British counterparts, and therefore would not fit British equipment properly. This, in part, helped lead to the development of numerous United States Military Standards and specifications for the manufacturing of essentially any piece of equipment that is used for military or defense purposes, including fasteners. World War II was a significant factor in this change.
A key component of most military standards is traceability. Put simply, hardware manufacturers must be able to trace their materials to their source, and provide traceability for their parts going into the supply chain, usually via bar codes or similar methods. This traceability is intended to help ensure that the right parts are used and that quality standards are met in each step of the manufacturing process; additionally, substandard parts can traced back to their source.
History
In 1988, the United States House Energy Subcommittee on Oversight and Investigations investigated counterfeit, mismarked, substandard fasteners and found extensive use in critical civilian and military infrastructure. As a result, they proposed Fastener Quality Assurance Act of 1988 (HR5051) that would require laboratory testing of fasteners in critical use applications prior to sale.
| Technology | Components_2 | null |
5024887 | https://en.wikipedia.org/wiki/Solar%20eclipse | Solar eclipse | A solar eclipse occurs when the Moon passes between Earth and the Sun, thereby obscuring the view of the Sun from a small part of Earth, totally or partially. Such an alignment occurs approximately every six months, during the eclipse season in its new moon phase, when the Moon's orbital plane is closest to the plane of Earth's orbit. In a total eclipse, the disk of the Sun is fully obscured by the Moon. In partial and annular eclipses, only part of the Sun is obscured. Unlike a lunar eclipse, which may be viewed from anywhere on the night side of Earth, a solar eclipse can only be viewed from a relatively small area of the world. As such, although total solar eclipses occur somewhere on Earth every 18 months on average, they recur at any given place only once every 360 to 410 years.
If the Moon were in a perfectly circular orbit and in the same orbital plane as Earth, there would be total solar eclipses once a month, at every new moon. Instead, because the Moon's orbit is tilted at about 5 degrees to Earth's orbit, its shadow usually misses Earth. Solar (and lunar) eclipses therefore happen only during eclipse seasons, resulting in at least two, and up to five, solar eclipses each year, no more than two of which can be total. Total eclipses are rarer because they require a more precise alignment between the centers of the Sun and Moon, and because the Moon's apparent size in the sky is sometimes too small to fully cover the Sun.
An eclipse is a natural phenomenon. In some ancient and modern cultures, solar eclipses were attributed to supernatural causes or regarded as bad omens. Astronomers' predictions of eclipses began in China as early as the 4th century BC; eclipses hundreds of years into the future may now be predicted with high accuracy.
Looking directly at the Sun can lead to permanent eye damage, so special eye protection or indirect viewing techniques are used when viewing a solar eclipse. Only the total phase of a total solar eclipse is safe to view without protection. Enthusiasts known as eclipse chasers or umbraphiles travel to remote locations to see solar eclipses.
Types
The Sun's distance from Earth is about 400 times the Moon's distance, and the Sun's diameter is about 400 times the Moon's diameter. Because these ratios are approximately the same, the Sun and the Moon as seen from Earth appear to be approximately the same size: about 0.5 degree of arc in angular measure.
The Moon's orbit around Earth is slightly elliptical, as is Earth's orbit around the Sun. The apparent sizes of the Sun and Moon therefore vary. The magnitude of an eclipse is the ratio of the apparent size of the Moon to the apparent size of the Sun during an eclipse. An eclipse that occurs when the Moon is near its closest distance to Earth (i.e., near its perigee) can be a total eclipse because the Moon will appear to be large enough to completely cover the Sun's bright disk or photosphere; a total eclipse has a magnitude greater than or equal to 1.000. Conversely, an eclipse that occurs when the Moon is near its farthest distance from Earth (i.e., near its apogee) can be only an annular eclipse because the Moon will appear to be slightly smaller than the Sun; the magnitude of an annular eclipse is less than 1.
Because Earth's orbit around the Sun is also elliptical, Earth's distance from the Sun similarly varies throughout the year. This affects the apparent size of the Sun in the same way, but not as much as does the Moon's varying distance from Earth. When Earth approaches its farthest distance from the Sun in early July, a total eclipse is somewhat more likely, whereas conditions favour an annular eclipse when Earth approaches its closest distance to the Sun in early January.
There are three main types of solar eclipses:
Total eclipse
A total eclipse occurs on average every 18 months when the dark silhouette of the Moon completely obscures the bright light of the Sun, allowing the much fainter solar corona to be visible. During an eclipse, totality occurs only along a narrow track on the surface of Earth. This narrow track is called the path of totality.
Annular eclipse
An annular eclipse, like a total eclipse, occurs when the Sun and Moon are exactly in line with Earth. During an annular eclipse, however, the apparent size of the Moon is not large enough to completely block out the Sun. Totality thus does not occur; the Sun instead appears as a very bright ring, or annulus, surrounding the dark disk of the Moon. Annular eclipses occur once every one or two years, not annually. The term derives from the Latin root word anulus, meaning "ring", rather than annus, for "year".
Partial eclipse
A partial eclipse occurs about twice a year, when the Sun and Moon are not exactly in line with Earth and the Moon only partially obscures the Sun. This phenomenon can usually be seen from a large part of Earth outside of the track of an annular or total eclipse. However, some eclipses can be seen only as a partial eclipse, because the umbra passes above Earth's polar regions and never intersects Earth's surface. Partial eclipses are virtually unnoticeable in terms of the Sun's brightness, as it takes well over 90% coverage to notice any darkening at all. Even at 99%, it would be no darker than civil twilight.
Terminology
Hybrid eclipse
A hybrid eclipse (also called annular/total eclipse) shifts between a total and annular eclipse. At certain points on the surface of Earth, it appears as a total eclipse, whereas at other points it appears as annular. Hybrid eclipses are comparatively rare.
A hybrid eclipse occurs when the magnitude of an eclipse changes during the event from less to greater than one, so the eclipse appears to be total at locations nearer the midpoint, and annular at other locations nearer the beginning and end, since the sides of Earth are slightly further away from the Moon. These eclipses are extremely narrow in their path width and relatively short in their duration at any point compared with fully total eclipses; the 2023 April 20 hybrid eclipse's totality is over a minute in duration at various points along the path of totality. Like a focal point, the width and duration of totality and annularity are near zero at the points where the changes between the two occur.
Central eclipse
Central eclipse is often used as a generic term for a total, annular, or hybrid eclipse. This is, however, not completely correct: the definition of a central eclipse is an eclipse during which the central line of the umbra touches Earth's surface. It is possible, though extremely rare, that part of the umbra intersects with Earth (thus creating an annular or total eclipse), but not its central line. This is then called a non-central total or annular eclipse. Gamma is a measure of how centrally the shadow strikes. The last (umbral yet) non-central solar eclipse was on April 29, 2014. This was an annular eclipse. The next non-central total solar eclipse will be on April 9, 2043.
Eclipse phases
The visual phases observed during a total eclipse are called:
First contact—when the Moon's limb (edge) is exactly tangential to the Sun's limb.
Second contact—starting with Baily's Beads (caused by light shining through valleys on the Moon's surface) and the diamond ring effect. Almost the entire disk is covered.
Totality—the Moon obscures the entire disk of the Sun and only the solar corona is visible.
Third contact—when the first bright light becomes visible and the Moon's shadow is moving away from the observer. Again a diamond ring may be observed.
Fourth contact—when the trailing edge of the Moon ceases to overlap with the solar disk and the eclipse ends.
Predictions
Geometry
The diagrams to the right show the alignment of the Sun, Moon, and Earth during a solar eclipse. The dark gray region between the Moon and Earth is the umbra, where the Sun is completely obscured by the Moon. The small area where the umbra touches Earth's surface is where a total eclipse can be seen. The larger light gray area is the penumbra, in which a partial eclipse can be seen. An observer in the antumbra, the area of shadow beyond the umbra, will see an annular eclipse.
The Moon's orbit around Earth is inclined at an angle of just over 5 degrees to the plane of Earth's orbit around the Sun (the ecliptic). Because of this, at the time of a new moon, the Moon will usually pass to the north or south of the Sun. A solar eclipse can occur only when a new moon occurs close to one of the points (known as nodes) where the Moon's orbit crosses the ecliptic.
As noted above, the Moon's orbit is also elliptical. The Moon's distance from Earth varies by up to about 5.9% from its average value. Therefore, the Moon's apparent size varies with its distance from Earth, and it is this effect that leads to the difference between total and annular eclipses. The distance of Earth from the Sun also varies during the year, but this is a smaller effect (by up to about 0.85% from its average value). On average, the Moon appears to be slightly (2.1%) smaller than the Sun as seen from Earth, so the majority (about 60%) of central eclipses are annular. It is only when the Moon is closer to Earth than average (near its perigee) that a total eclipse occurs.
The Moon orbits Earth in approximately 27.3 days, relative to a fixed frame of reference. This is known as the sidereal month. However, during one sidereal month, Earth has revolved part way around the Sun, making the average time between one new moon and the next longer than the sidereal month: it is approximately 29.5 days. This is known as the synodic month and corresponds to what is commonly called the lunar month.
The Moon crosses from south to north of the ecliptic at its ascending node, and vice versa at its descending node. However, the nodes of the Moon's orbit are gradually moving in a retrograde motion, due to the action of the Sun's gravity on the Moon's motion, and they make a complete circuit every 18.6 years. This regression means that the time between each passage of the Moon through the ascending node is slightly shorter than the sidereal month. This period is called the nodical or draconic month.
Finally, the Moon's perigee is moving forwards or precessing in its orbit and makes a complete circuit in 8.85 years. The time between one perigee and the next is slightly longer than the sidereal month and known as the anomalistic month.
The Moon's orbit intersects with the ecliptic at the two nodes that are 180 degrees apart. Therefore, the new moon occurs close to the nodes at two periods of the year approximately six months (173.3 days) apart, known as eclipse seasons, and there will always be at least one solar eclipse during these periods. Sometimes the new moon occurs close enough to a node during two consecutive months to eclipse the Sun on both occasions in two partial eclipses. This means that, in any given year, there will always be at least two solar eclipses, and there can be as many as five.
Eclipses can occur only when the Sun is within about 15 to 18 degrees of a node, (10 to 12 degrees for central eclipses). This is referred to as an eclipse limit, and is given in ranges because the apparent sizes and speeds of the Sun and Moon vary throughout the year. In the time it takes for the Moon to return to a node (draconic month), the apparent position of the Sun has moved about 29 degrees, relative to the nodes. Since the eclipse limit creates a window of opportunity of up to 36 degrees (24 degrees for central eclipses), it is possible for partial eclipses (or rarely a partial and a central eclipse) to occur in consecutive months.
Path
During a central eclipse, the Moon's umbra (or antumbra, in the case of an annular eclipse) moves rapidly from west to east across Earth. Earth is also rotating from west to east, at about 28 km/min at the Equator, but as the Moon is moving in the same direction as Earth's rotation at about 61 km/min, the umbra almost always appears to move in a roughly west–east direction across a map of Earth at the speed of the Moon's orbital velocity minus Earth's rotational velocity.
The width of the track of a central eclipse varies according to the relative apparent diameters of the Sun and Moon. In the most favourable circumstances, when a total eclipse occurs very close to perigee, the track can be up to wide and the duration of totality may be over 7 minutes. Outside of the central track, a partial eclipse is seen over a much larger area of Earth. Typically, the umbra is 100–160 km wide, while the penumbral diameter is in excess of 6400 km.
Besselian elements are used to predict whether an eclipse will be partial, annular, or total (or annular/total), and what the eclipse circumstances will be at any given location.
Calculations with Besselian elements can determine the exact shape of the umbra's shadow on Earth's surface. But at what longitudes on Earth's surface the shadow will fall, is a function of Earth's rotation, and on how much that rotation has slowed down over time. A number called ΔT is used in eclipse prediction to take this slowing into account. As Earth slows, ΔT increases. ΔT for dates in the future can only be roughly estimated because Earth's rotation is slowing irregularly. This means that, although it is possible to predict that there will be a total eclipse on a certain date in the far future, it is not possible to predict in the far future exactly at what longitudes that eclipse will be total. Historical records of eclipses allow estimates of past values of ΔT and so of Earth's rotation.
Duration
The following factors determine the duration of a total solar eclipse (in order of decreasing importance):
The Moon being almost exactly at perigee (making its angular diameter as large as possible).
Earth being very near aphelion (furthest away from the Sun in its elliptical orbit, making its angular diameter nearly as small as possible).
The midpoint of the eclipse being very close to Earth's equator, where the rotational velocity is greatest and is closest to the speed of the lunar shadow moving over Earth's surface.
The vector of the eclipse path at the midpoint of the eclipse aligning with the vector of Earth's rotation (i.e. not diagonal but due east).
The midpoint of the eclipse being near the subsolar point (the part of Earth closest to the Sun).
The longest eclipse that has been calculated thus far is the eclipse of July 16, 2186 (with a maximum duration of 7 minutes 29 seconds over northern Guyana).
Occurrence and cycles
A total solar eclipse is a rare event, recurring somewhere on Earth every 18 months on average, yet is estimated to recur at any given location only every 360–410 years on average. The total eclipse lasts for only a maximum of a few minutes at any location because the Moon's umbra moves eastward at over . Totality currently can never last more than 7 min 32 s. This value changes over the millennia and is currently decreasing. By the 8th millennium, the longest theoretically possible total eclipse will be less than 7 min 2 s. The last time an eclipse longer than 7 minutes occurred was June 30, 1973 (7 min 3 sec). Observers aboard a Concorde supersonic aircraft were able to stretch totality for this eclipse to about 74 minutes by flying along the path of the Moon's umbra. The next total eclipse exceeding seven minutes in duration will not occur until June 25, 2150. The longest total solar eclipse during the year period from 3000 BC to at least 8000 AD will occur on July 16, 2186, when totality will last 7 min 29 s. For comparison, the longest total eclipse of the 20th century at 7 min 8 s occurred on June 20, 1955, and there will be no total solar eclipses over 7 min in duration in the 21st century.
It is possible to predict other eclipses using eclipse cycles. The saros is probably the best known and one of the most accurate. A saros lasts 6585.3 days (a little over 18 years), which means that, after this period, a practically identical eclipse will occur. The most notable difference will be a westward shift of about 120° in longitude (due to the 0.3 days) and a little in latitude (north-south for odd-numbered cycles, the reverse for even-numbered ones). A saros series always starts with a partial eclipse near one of Earth's polar regions, then shifts over the globe through a series of annular or total eclipses, and ends with a partial eclipse at the opposite polar region. A saros series lasts 1226 to 1550 years and 69 to 87 eclipses, with about 40 to 60 of them being central.
Frequency per year
Between two and five solar eclipses occur every year, with at least one per eclipse season. Since the Gregorian calendar was instituted in 1582, years that have had five solar eclipses were 1693, 1758, 1805, 1823, 1870, and 1935. The next occurrence will be 2206. On average, there are about 240 solar eclipses each century.
Final totality
Total solar eclipses are seen on Earth because of a fortuitous combination of circumstances. Even on Earth, the diversity of eclipses familiar to people today is a temporary (on a geological time scale) phenomenon. Hundreds of millions of years in the past, the Moon was closer to Earth and therefore apparently larger, so every solar eclipse was total or partial, and there were no annular eclipses. Due to tidal acceleration, the orbit of the Moon around Earth becomes approximately 3.8 cm more distant each year. Millions of years in the future, the Moon will be too far away to fully occlude the Sun, and no total eclipses will occur. In the same timeframe, the Sun may become brighter, making it appear larger in size. Estimates of the time when the Moon will be unable to occlude the entire Sun when viewed from Earth range between 650 million and 1.4 billion years in the future.
Viewing
Looking directly at the photosphere of the Sun (the bright disk of the Sun itself), even for just a few seconds, can cause permanent damage to the retina of the eye, because of the intense visible and invisible radiation that the photosphere emits. This damage can result in impairment of vision, up to and including blindness. The retina has no sensitivity to pain, and the effects of retinal damage may not appear for hours, so there is no warning that injury is occurring.
Under normal conditions, the Sun is so bright that it is difficult to stare at it directly. However, during an eclipse, with so much of the Sun covered, it is easier and more tempting to stare at it. Looking at the Sun during an eclipse is as dangerous as looking at it outside an eclipse, except during the brief period of totality, when the Sun's disk is completely covered (totality occurs only during a total eclipse and only very briefly; it does not occur during a partial or annular eclipse). Viewing the Sun's disk through any kind of optical aid (binoculars, a telescope, or even an optical camera viewfinder) is extremely hazardous and can cause irreversible eye damage within a fraction of a second.
Partial and annular eclipses
Viewing the Sun during partial and annular eclipses (and during total eclipses outside the brief period of totality) requires special eye protection, or indirect viewing methods if eye damage is to be avoided. The Sun's disk can be viewed using appropriate filtration to block the harmful part of the Sun's radiation. Sunglasses do not make viewing the Sun safe. Only properly designed and certified solar filters should be used for direct viewing of the Sun's disk. Especially, self-made filters using common objects such as a floppy disk removed from its case, a Compact Disc, a black colour slide film, smoked glass, etc. must be avoided.
The safest way to view the Sun's disk is by indirect projection. This can be done by projecting an image of the disk onto a white piece of paper or card using a pair of binoculars (with one of the lenses covered), a telescope, or another piece of cardboard with a small hole in it (about 1 mm diameter), often called a pinhole camera. The projected image of the Sun can then be safely viewed; this technique can be used to observe sunspots, as well as eclipses. Care must be taken, however, to ensure that no one looks through the projector (telescope, pinhole, etc.) directly. A kitchen colander with small holes can also be used to project multiple images of the partially eclipsed Sun onto the ground or a viewing screen. Viewing the Sun's disk on a video display screen (provided by a video camera or digital camera) is safe, although the camera itself may be damaged by direct exposure to the Sun. The optical viewfinders provided with some video and digital cameras are not safe. Securely mounting #14 welder's glass in front of the lens and viewfinder protects the equipment and makes viewing possible. Professional workmanship is essential because of the dire consequences any gaps or detaching mountings will have. In the partial eclipse path, one will not be able to see the corona or nearly complete darkening of the sky. However, depending on how much of the Sun's disk is obscured, some darkening may be noticeable. If three-quarters or more of the Sun is obscured, then an effect can be observed by which the daylight appears to be dim, as if the sky were overcast, yet objects still cast sharp shadows.
Totality
When the shrinking visible part of the photosphere becomes very small, Baily's beads will occur. These are caused by the sunlight still being able to reach Earth through lunar valleys. Totality then begins with the diamond ring effect, the last bright flash of sunlight.
It is safe to observe the total phase of a solar eclipse directly only when the Sun's photosphere is completely covered by the Moon, and not before or after totality. During this period, the Sun is too dim to be seen through filters. The Sun's faint corona will be visible, and the chromosphere, solar prominences, coronal streamers and possibly even a solar flare may be seen. At the end of totality, the same effects will occur in reverse order, and on the opposite side of the Moon.
Eclipse chasing
A dedicated group of eclipse chasers have pursued the observation of solar eclipses when they occur around Earth. A person who chases eclipses is known as an umbraphile, meaning shadow lover. Umbraphiles travel for eclipses and use various tools to help view the sun including solar viewing glasses, also known as eclipse glasses, as well as telescopes.
Photography
The first known photograph of a solar eclipse was taken on July 28, 1851, by Johann Julius Friedrich Berkowski, using the daguerreotype process.
Photographing an eclipse is possible with fairly common camera equipment. In order for the disk of the Sun/Moon to be easily visible, a fairly high magnification long focus lens is needed (at least 200 mm for a 35 mm camera), and for the disk to fill most of the frame, a longer lens is needed (over 500 mm). As with viewing the Sun directly, looking at it through the optical viewfinder of a camera can produce damage to the retina, so care is recommended. Solar filters are required for digital photography even if an optical viewfinder is not used. Using a camera's live view feature or an electronic viewfinder is safe for the human eye, but the Sun's rays could potentially irreparably damage digital image sensors unless the lens is covered by a properly designed solar filter.
Historical eclipses
Historical eclipses are a very valuable resource for historians, in that they allow a few historical events to be dated precisely, from which other dates and ancient calendars may be deduced. The oldest recorded solar eclipse was recorded on a clay tablet found at Ugarit, in modern Syria, with two plausible dates usually cited: 3 May 1375 BC or 5 March 1223 BC, the latter being favored by most recent authors on the topic. A solar eclipse of June 15, 763 BC mentioned in an Assyrian text is important for the chronology of the ancient Near East. There have been other claims to date earlier eclipses. The legendary Chinese king Zhong Kang supposedly beheaded two astronomers, Hsi and Ho, who failed to predict an eclipse 4000 years ago. Perhaps the earliest still-unproven claim is that of archaeologist Bruce Masse, who putatively links an eclipse that occurred on May 10, 2807, BC with a possible meteor impact in the Indian Ocean on the basis of several ancient flood myths that mention a total solar eclipse.
Eclipses have been interpreted as omens, or portents. The ancient Greek historian Herodotus wrote that Thales of Miletus predicted an eclipse that occurred during a battle between the Medes and the Lydians. Both sides put down their weapons and declared peace as a result of the eclipse. The exact eclipse involved remains uncertain, although the issue has been studied by hundreds of ancient and modern authorities. One likely candidate took place on May 28, 585 BC, probably near the Halys river in Asia Minor. An eclipse recorded by Herodotus before Xerxes departed for his expedition against Greece, which is traditionally dated to 480 BC, was matched by John Russell Hind to an annular eclipse of the Sun at Sardis on February 17, 478 BC. Alternatively, a partial eclipse was visible from Persia on October 2, 480 BC. Herodotus also reports a solar eclipse at Sparta during the Second Persian invasion of Greece. The date of the eclipse (August 1, 477 BC) does not match exactly the conventional dates for the invasion accepted by historians.
In ancient China, where solar eclipses were known as an "eating of the Sun" ( ), the earliest records of eclipses date to around 720 BC. The 4th century BC astronomer Shi Shen described the prediction of eclipses by using the relative positions of the Moon and Sun.
Attempts have been made to establish the exact date of Good Friday by assuming that the darkness described at Jesus's crucifixion was a solar eclipse. This research has not yielded conclusive results, and Good Friday is recorded as being at Passover, which is held at the time of a full moon. Further, the darkness lasted from the sixth hour to the ninth, or three hours, which is much, much longer than the eight-minute upper limit for any solar eclipse's totality. Contemporary chronicles wrote about an eclipse at the beginning of May 664 that coincided with the beginning of the plague of 664 in the British isles. In the Western hemisphere, there are few reliable records of eclipses before AD 800, until the advent of Arab and monastic observations in the early medieval period.
A solar eclipse took place on January 27, 632 over Arabia during Muhammad's lifetime. Muhammad denied the eclipse had anything to do with his son dying earlier that day, saying "The sun and the moon do not eclipse because of the death of someone from the people but they are two signs amongst the signs of God." The Cairo astronomer Ibn Yunus wrote that the calculation of eclipses was one of the many things that connect astronomy with the Islamic law, because it allowed knowing when a special prayer can be made. The first recorded observation of the corona was made in Constantinople in AD 968.
The first known telescopic observation of a total solar eclipse was made in France in 1706. Nine years later, English astronomer Edmund Halley accurately predicted and observed the solar eclipse of May 3, 1715. By the mid-19th century, scientific understanding of the Sun was improving through observations of the Sun's corona during solar eclipses. The corona was identified as part of the Sun's atmosphere in 1842, and the first photograph (or daguerreotype) of a total eclipse was taken of the solar eclipse of July 28, 1851. Spectroscope observations were made of the solar eclipse of August 18, 1868, which helped to determine the chemical composition of the Sun.
John Fiske summed up myths about the solar eclipse like this in his 1872 book Myth and Myth-Makers,
Particular observations, phenomena and impact
A total solar eclipse provides a rare opportunity to observe the corona (the outer layer of the Sun's atmosphere). Normally this is not visible because the photosphere is much brighter than the corona. According to the point reached in the solar cycle, the corona may appear small and symmetric, or large and fuzzy. It is very hard to predict this in advance.
Phenomena associated with eclipses include shadow bands (also known as flying shadows), which are similar to shadows on the bottom of a swimming pool. They occur only just prior to and after totality, when a narrow solar crescent acts as an anisotropic light source. As the light filters through leaves of trees during a partial eclipse, the overlapping leaves create natural pinholes, displaying mini eclipses on the ground.
1919 observations
The observation of a total solar eclipse of May 29, 1919, helped to confirm Einstein's theory of general relativity. By comparing the apparent distance between stars in the constellation Taurus, with and without the Sun between them, Arthur Eddington stated that the theoretical predictions about gravitational lenses were confirmed. The observation with the Sun between the stars was possible only during totality since the stars are then visible. Though Eddington's observations were near the experimental limits of accuracy at the time, work in the later half of the 20th century confirmed his results.
Gravity anomalies
There is a long history of observations of gravity-related phenomena during solar eclipses, especially during the period of totality. Maurice Allais reported observing unusual and unexplained movements during solar eclipses in both 1954 and 1959. The reality of this phenomenon, named the Allais effect, has remained controversial. Similarly, in 1970, Saxl and Allen observed the sudden change in motion of a torsion pendulum; this phenomenon is called the Saxl effect.
Observation during the 1997 solar eclipse by Wang et al. suggested a possible gravitational shielding effect, which generated debate. In 2002, Wang and a collaborator published detailed data analysis, which suggested that the phenomenon still remains unexplained.
Eclipses and transits
In principle, the simultaneous occurrence of a solar eclipse and a transit of a planet is possible. But these events are extremely rare because of their short durations. The next anticipated simultaneous occurrence of a solar eclipse and a transit of Mercury will be on July 5, 6757, and a solar eclipse and a transit of Venus is expected on April 5, .
More common, but still infrequent, is a conjunction of a planet (especially, but not only, Mercury or Venus) at the time of a total solar eclipse, in which event the planet will be visible very near the eclipsed Sun, when without the eclipse it would have been lost in the Sun's glare. At one time, some scientists hypothesized that there may be a planet (often given the name Vulcan) even closer to the Sun than Mercury; the only way to confirm its existence would have been to observe it in transit or during a total solar eclipse. No such planet was ever found, and general relativity has since explained the observations that led astronomers to suggest that Vulcan might exist.
Artificial satellites
Artificial satellites can also pass in front of the Sun as seen from Earth, but none is large enough to cause an eclipse. At the altitude of the International Space Station, for example, an object would need to be about across to blot the Sun out entirely. These transits are difficult to watch because the zone of visibility is very small. The satellite passes over the face of the Sun in about a second, typically. As with a transit of a planet, it will not get dark.
Observations of eclipses from spacecraft or artificial satellites orbiting above Earth's atmosphere are not subject to weather conditions. The crew of Gemini 12 observed a total solar eclipse from space in 1966. The partial phase of the 1999 total eclipse was visible from Mir.
Impact
The solar eclipse of March 20, 2015, was the first occurrence of an eclipse estimated to potentially have a significant impact on the power system, with the electricity sector taking measures to mitigate any impact. The continental Europe and Great Britain synchronous areas were estimated to have about 90 gigawatts of solar power and it was estimated that production would temporarily decrease by up to 34 GW compared to a clear sky day.
Eclipses may cause the temperature to decrease by , with wind power potentially decreasing as winds are reduced by per second.
In addition to the drop in light level and air temperature, animals change their behavior during totality. For example, birds and squirrels return to their nests and crickets chirp.
Recent and forthcoming solar eclipses
Eclipses occur only in the eclipse season, when the Sun is close to either the ascending or descending node of the Moon. Each eclipse is separated by one, five or six lunations (synodic months), and the midpoint of each season is separated by 173.3 days, which is the mean time for the Sun to travel from one node to the next. The period is a little less than half a calendar year because the lunar nodes slowly regress. Because 223 synodic months is roughly equal to 239 anomalistic months and 242 draconic months, eclipses with similar geometry recur 223 synodic months (about 6,585.3 days) apart. This period (18 years 11.3 days) is a saros. Because 223 synodic months is not identical to 239 anomalistic months or 242 draconic months, saros cycles do not endlessly repeat. Each cycle begins with the Moon's shadow crossing Earth near the north or south pole, and subsequent events progress toward the other pole until the Moon's shadow misses Earth and the series ends. Saros cycles are numbered; currently, cycles 117 to 156 are active.
2018–2021
2022–2025
2026–2029
| Physical sciences | Celestial mechanics | null |
3702612 | https://en.wikipedia.org/wiki/True%20toad | True toad | A true toad is any member of the family Bufonidae, in the order Anura (frogs and toads). This is the only family of anurans in which all members are known as toads, although some may be called frogs (such as harlequin frogs). The bufonids now comprise more than 35 genera, Bufo being the best known.
History
Bufonidae is thought to have originated in South America. Some studies date the origin of the group to after the breakup of Gondwana, about 78–99 million years ago in the Late Cretaceous. In contrast, other studies have dated the origin of the group to the early Paleocene. The bufonids likely radiated out of South America during the Eocene, with the entire radiation occurring during the Eocene to Oligocene, marking an extremely rapid divergence likely facilitated by the Paleogene's changing climatic conditions.
Taxonomy
The following phylogeny of most genera in the family is based on Portik and Papenfuss, 2015:, Chan et al., 2016, Chandramouli et al., 2016, and Kok et al., 2017
Ingerophrynus alongside Leptophryne was grouped as basal to the clade containing all other Southeast Asian toad genera and Ghatophryne by Portik and Papenfuss, but was found to group with Phrynoidis and Rentapia by Chan et al. Ghatophryne was grouped with Phrynoidis and Rentapia by Portik and Papenfuss but was found to group with Pelophryne and Ansonia by Chan et al. In addition, Sabahphrynus was grouped with Strauchbufo and Bufo by Portik and Papenfuss but was found to group with Pelophryne, Ansonia, and Ghatophryne by Chan et al.
Characteristics
True toads are widespread and are native to every continent except Australia and Antarctica, inhabiting a variety of environments, from arid areas to rainforest. Most lay eggs in paired strings that hatch into tadpoles, although, in the genus Nectophrynoides, the eggs hatch directly into miniature toads.
All true toads are toothless and generally warty in appearance. They have a pair of parotoid glands on the back of their heads. These glands contain an alkaloid poison which the toads excrete when stressed. The poison in the glands contains a number of toxins causing different effects. Bufotoxin is a general term. Different animals contain significantly different substances and proportions of substances. Some, like the cane toad Rhinella marina, are more toxic than others. Some "psychoactive toads", such as the Colorado River toad Incilius alvarius, have been used recreationally for the effects of their bufotoxin.
Depending on the species, male or female toads may possess a Bidder's organ, a trait unique to all bufonids except genera Melanophryniscus and Truebella. Under the right conditions, the organ becomes an active ovary.
The loss of teeth has arisen in frogs independently over 20 times. Notably, all members of Bufonidae are toothless. Another Anuran family with a comparable degree of edentulism is the family Microhylidae.
Reproduction
Internal fertilization occurs in four bufonid genera.
Mertensophryne (some species)
Nectophrynoides (presumably all species)
Altiphrynoides malcolmi (one out of two species in the genus Altiphrynoides)
Nimbaphrynoides occidentalis (the sole species in the monotypic genus Nimbaphrynoides)
Ascaphus (all species) and Eleutherodactylus (two species, E. coqui and E. jasperi) are the only other frog genera that have internal fertilization. Limnonectes larvaepartus also has internal fertilization.
Taxonomy and genera
The family Bufonidae contains over 570 species among 52 genera.
The family also contains an incertae sedis species, "Bufo" scorteccii .
| Biology and health sciences | Amphibians | null |
3704228 | https://en.wikipedia.org/wiki/Pipe%20%28fluid%20conveyance%29 | Pipe (fluid conveyance) | A pipe is a tubular section or hollow cylinder, usually but not necessarily of circular cross-section, used mainly to convey substances which can flow — liquids and gases (fluids), slurries, powders and masses of small solids. It can also be used for structural applications; a hollow pipe is far stiffer per unit weight than the solid members.
In common usage the words pipe and tube are usually interchangeable, but in industry and engineering, the terms are uniquely defined. Depending on the applicable standard to which it is manufactured, pipe is generally specified by a nominal diameter with a constant outside diameter (OD) and a schedule that defines the thickness. Tube is most often specified by the OD and wall thickness, but may be specified by any two of OD, inside diameter (ID), and wall thickness. Pipe is generally manufactured to one of several international and national industrial standards. While similar standards exist for specific industry application tubing, tube is often made to custom sizes and a broader range of diameters and tolerances. Many industrial and government standards exist for the production of pipe and tubing. The term "tube" is also commonly applied to non-cylindrical sections, i.e., square or rectangular tubing. In general, "pipe" is the more common term in most of the world, whereas "tube" is more widely used in the United States.
Both "pipe" and "tube" imply a level of rigidity and permanence, whereas a hose (or hosepipe) is usually portable and flexible. Pipe assemblies are almost always constructed with the use of fittings such as elbows, tees, and so on, while tube may be formed or bent into custom configurations. For materials that are inflexible, cannot be formed, or where construction is governed by codes or standards, tube assemblies are also constructed with the use of tube fittings.
Uses
Plumbing
Tap water
Irrigation
Pipelines transporting gas or liquid over long distances
Compressed air systems
Casing for concrete pilings used in construction projects
High-temperature or high-pressure manufacturing processes
The petroleum industry:
Oil well casing
Oil refinery equipment
Delivery of fluids, either gaseous or liquid, in a process plant from one point to another point in the process
Delivery of bulk solids, in a food or process plant from one point to another point in the process
The construction of high pressure storage vessels (large pressure vessels are constructed from plate, not pipe owing to their wall thickness and size).
Additionally, pipes are used for many purposes that do not involve conveying fluid. Handrails, scaffolding, and support structures are often constructed from structural pipes, especially in an industrial environment.
History
The first known use of pipes was in Ancient Egypt. The Pyramid of Sahure, completed around the 25th century BC, included a temple with an elaborate drainage system including more than of copper piping.
During the Napoleonic Wars Birmingham gunmakers tried to use rolling mills to make iron musket barrels. One of them, Henry Osborne, developed a relatively effective process in 1817 with which he started to make iron gas tubes ca. 1820, selling some to gas lighting pioneer Samuel Clegg.
When steel pipes were introduced in 19th century, they initially were riveted, and later clamped with H-shaped bars (even though methods for making weldless steel tubes were known already in the 1870s), until by the early 1930s these methods were replaced by welding, which is still widely used today.
Manufacture
There are three processes for metallic pipe manufacture. Centrifugal casting of hot alloyed metal is one of the most prominent process. Ductile iron pipes are generally manufactured in such a fashion.
Seamless pipe (SMLS) is formed by drawing a solid billet over a piercing rod to create the hollow shell in a process called rotary piercing. As the manufacturing process does not include any welding, seamless pipes are perceived to be stronger and more reliable. Historically, seamless pipe was regarded as withstanding pressure better than other types, and was often more available than welded pipe.
Advances since the 1970s, in materials, process control, and non-destructive testing, allow correctly specified welded pipe to replace seamless in many applications. Welded pipe is formed by rolling plate and welding the seam (usually by Electric resistance welding ("ERW"), or Electric Fusion Welding ("EFW")). The weld flash can be removed from both inner and outer surfaces using a scarfing blade. The weld zone can also be heat-treated to make the seam less visible. Welded pipe often has tighter dimensional tolerances than the seamless type, and can be cheaper to manufacture.
There are a number of processes that may be used to produce ERW pipes. Each of these processes leads to coalescence or merging of steel components into pipes. Electric current is passed through the surfaces that have to be welded together; as the components being welded together resist the electric current, heat is generated which forms the weld. Pools of molten metal are formed where the two surfaces are connected as a strong electric current is passed through the metal; these pools of molten metal form the weld that binds the two abutted components.
ERW pipes are manufactured from the longitudinal welding of steel. The welding process for ERW pipes is continuous, as opposed to welding of distinct sections at intervals. ERW process uses steel coil as feedstock.
The High Frequency Induction Technology (HFI) welding process is used for manufacturing ERW pipes. In this process, the current to weld the pipe is applied by means of an induction coil around the tube. HFI is generally considered to be technically superior to "ordinary" ERW when manufacturing pipes for critical applications, such as for usage in the energy sector, in addition to other uses in line pipe applications, as well as for casing and tubing.
Large-diameter pipe ( or greater) may be ERW, EFW, or Submerged Arc Welded ("SAW") pipe. There are two technologies that can be used to manufacture steel pipes of sizes larger than the steel pipes that can be produced by seamless and ERW processes. The two types of pipes produced through these technologies are longitudinal-submerged arc-welded (LSAW) and spiral-submerged arc-welded (SSAW) pipes. LSAW are made by bending and welding wide steel plates and most commonly used in oil and gas industry applications. Due to their high cost, LSAW pipes are seldom used in lower value non-energy applications such as water pipelines. SSAW pipes are produced by spiral (helicoidal) welding of steel coil and have a cost advantage over LSAW pipes, as the process uses coils rather than steel plates. As such, in applications where spiral-weld is acceptable, SSAW pipes may be preferred over LSAW pipes. Both LSAW pipes and SSAW pipes compete against ERW pipes and seamless pipes in the diameter ranges of 16”-24”.
Tubing for flow, either metal or plastic, is generally extruded.
Materials
Pipe is made out of many types of material including ceramic, glass, fiberglass, many metals, concrete and plastic. In the past, wood and lead (Latin plumbum, from which comes the word 'plumbing') were commonly used.
Typically metallic piping is made of steel or iron, such as unfinished, black (lacquer) steel, carbon steel, stainless steel, galvanized steel, brass, and ductile iron. Iron based piping is subject to corrosion if used within a highly oxygenated water stream. Aluminum pipe or tubing may be utilized where iron is incompatible with the service fluid or where weight is a concern; aluminum is also used for heat transfer tubing such as in refrigerant systems. Copper tubing is popular for domestic water (potable) plumbing systems; copper may be used where heat transfer is desirable (i.e. radiators or heat exchangers). Inconel, chrome moly, and titanium steel alloys are used in high temperature and pressure piping in process and power facilities. When specifying alloys for new processes, the known issues of creep and sensitization effect must be taken into account.
Lead piping is still found in old domestic and other water distribution systems, but is no longer permitted for new potable water piping installations due to its toxicity. Many building codes now require that lead piping in residential or institutional installations be replaced with non-toxic piping or that the tubes' interiors be treated with phosphoric acid. According to a senior researcher and lead expert with the Canadian Environmental Law Association, "[...] there is no safe level of lead [for human exposure]". In 1991 the US EPA issued the Lead and Copper Rule, a federal regulation which limits the concentration of lead and copper allowed in public drinking water, as well as the permissible amount of pipe corrosion occurring due to the water itself. In the US it is estimated that 6.5 million lead service lines (pipes that connect water mains to home plumbing) installed before the 1930s are still in use.
Plastic tubing is widely used for its light weight, chemical resistance, non-corrosive properties, and ease of making connections. Plastic materials include polyvinyl chloride (PVC), chlorinated polyvinyl chloride (CPVC), fibre reinforced plastic (FRP), reinforced polymer mortar (RPMP), polypropylene (PP), polyethylene (PE), cross-linked high-density polyethylene (PEX), polybutylene (PB), and acrylonitrile butadiene styrene (ABS), for example. In many countries, PVC pipes account for most pipe materials used in buried municipal applications for drinking water distribution and wastewater mains.
Pipe may be made from concrete or ceramic, usually for low-pressure applications such as gravity flow or drainage. Pipes for sewage are still predominantly made from concrete or vitrified clay. Reinforced concrete can be used for large-diameter concrete pipes. This pipe material can be used in many types of construction, and is often used in the gravity-flow transport of storm water. Usually such pipe will have a receiving bell or a stepped fitting, with various sealing methods applied at installation.
Traceability and positive material identification (PMI)
When the alloys for piping are forged, metallurgical tests are performed to determine material composition by % of each chemical element in the piping, and the results are recorded in a material test report, also known as a Mill Test Report (MTR). These tests can be used to prove that the alloy conforms to various specifications (e.g. 316 SS). The tests are stamped by the mill's QA/QC department and can be used to trace the material back to the mill by future users, such as piping and fitting manufacturers. Maintaining the traceability between the alloy material and associated MTR is an important quality assurance issue. QA often requires the heat number to be written on the pipe. Precautions must also be taken to prevent the introduction of counterfeit materials. As a backup to etching/labeling of the material identification on the pipe, positive material identification (PMI) is performed using a handheld device; the device scans the pipe material using an emitted electromagnetic wave (x-ray fluorescence/XRF) and receives a reply that is spectrographically analyzed.
Sizes
Pipe sizes can be confusing because the terminology may relate to historical dimensions. For example, a half-inch iron pipe does not have any dimension that is a half inch. Initially, a half inch pipe did have an inner diameter of —but it also had thick walls. As technology improved, thinner walls became possible, but the outside diameter stayed the same so it could mate with existing older pipe, increasing the inner diameter beyond half an inch. The history of copper pipe is similar. In the 1930s, the pipe was designated by its internal diameter and a wall thickness. Consequently, a copper pipe had a outside diameter. The outside diameter was the important dimension for mating with fittings. The wall thickness on modern copper is usually thinner than , so the internal diameter is only "nominal" rather than a controlling dimension. Newer pipe technologies sometimes adopted a sizing system as its own. PVC pipe uses the Nominal Pipe Size.
Pipe sizes are specified by a number of national and international standards, including API 5L, ANSI/ASME B36.10M and B36.19M in the US, BS 1600 and BS EN 10255 in the United Kingdom and Europe.
There are two common methods for designating pipe outside diameter (OD). The North American method is called NPS ("Nominal Pipe Size") and is based on inches (also frequently referred to as NB ("Nominal Bore")). The European version is called DN ("Diametre Nominal" / "Nominal Diameter") and is based on millimetres. Designating the outside diameter allows pipes of the same size to be fit together no matter what the wall thickness.
For pipe sizes less than NPS 14 inch (DN 350), both methods give a nominal value for the OD that is rounded off and is not the same as the actual OD. For example, NPS 2 inch and DN 50 are the same pipe, but the actual OD is . The only way to obtain the actual OD is to look it up in a reference table.
For pipe sizes of NPS 14 inch (DN 350) and greater the NPS size is the actual diameter in inches and the DN size is equal to NPS times 25 (not 25.4) rounded to a convenient multiple of 50. For example, NPS 14 has an OD of , and is equivalent to DN 350.
Since the outside diameter is fixed for a given pipe size, the inside diameter will vary depending on the wall thickness of the pipe. For example, 2" Schedule 80 pipe has thicker walls and therefore a smaller inside diameter than 2" Schedule 40 pipe.
Steel pipe has been produced for about 150 years. The pipe sizes that are in use today in PVC and galvanized were originally designed years ago for steel pipe. The number system, like Sch 40, 80, 160, were set long ago and seem a little odd. For example, Sch 20 pipe is even thinner than Sch 40, but same OD. And while these pipes are based on old steel pipe sizes, there is other pipe, like cpvc for heated water, that uses pipe sizes, inside and out, based on old copper pipe size standards instead of steel.
Many different standards exist for pipe sizes, and their prevalence varies depending on industry and geographical area. The pipe size designation generally includes two numbers; one that indicates the outside (OD) or nominal diameter, and the other that indicates the wall thickness. In the early twentieth century, American pipe was sized by inside diameter. This practice was abandoned to improve compatibility with pipe fittings that must usually fit the OD of the pipe, but it has had a lasting impact on modern standards around the world.
In North America and the UK, pressure piping is usually specified by Nominal Pipe Size (NPS) and schedule (SCH). Pipe sizes are documented by a number of standards, including API 5L, ANSI/ASME B36.10M (Table 1) in the US, and BS 1600 and BS 1387 in the United Kingdom. Typically the pipe wall thickness is the controlled variable, and the Inside Diameter (I.D.) is allowed to vary. The pipe wall thickness has a variance of approximately 12.5 percent.
In the rest of Europe pressure piping uses the same pipe IDs and wall thicknesses as Nominal Pipe Size, but labels them with a metric Diameter Nominal (DN) instead of the imperial NPS. For NPS larger than 14, the DN is equal to the NPS multiplied by 25. (Not 25.4) This is documented by EN 10255 (formerly DIN 2448 and BS 1387) and ISO 65:1981, and it is often called DIN or ISO pipe.
Japan has its own set of standard pipe sizes, often called JIS pipe.
The Iron pipe size (IPS) is an older system still used by some manufacturers and legacy drawings and equipment. The IPS number is the same as the NPS number, but the schedules were limited to Standard Wall (STD), Extra Strong (XS), and Double Extra Strong (XXS). STD is identical to SCH 40 for NPS 1/8 to NPS 10, inclusive, and indicates .375" wall thickness for NPS 12 and larger. XS is identical to SCH 80 for NPS 1/8 to NPS 8, inclusive, and indicates .500" wall thickness for NPS 8 and larger. Different definitions exist for XXS, however it is never the same as SCH 160. XXS is in fact thicker than SCH 160 for NPS 1/8" to 6" inclusive, whereas SCH 160 is thicker than XXS for NPS 8" and larger.
Another old system is the Ductile Iron Pipe Size (DIPS), which generally has larger ODs than IPS.
Copper plumbing tube for residential plumbing follows an entirely different size system in America, often called Copper Tube Size (CTS); see domestic water system. Its nominal size is neither the inside nor outside diameter. Plastic tubing, such as PVC and CPVC, for plumbing applications also has different sizing standards.
Agricultural applications use PIP sizes, which stands for Plastic Irrigation Pipe. PIP comes in pressure ratings of , , , , and and is generally available in diameters of .
Standards
The manufacture and installation of pressure piping is tightly regulated by the ASME "B31" code series such as B31.1 or B31.3 which have their basis in the ASME Boiler and Pressure Vessel Code (BPVC). This code has the force of law in Canada and the US. Europe and the rest of the world has an equivalent system of codes. Pressure piping is generally pipe that must carry pressures greater than 10 to 25 atmospheres, although definitions vary. To ensure safe operation of the system, the manufacture, storage, welding, testing, etc. of pressure piping must meet stringent quality standards.
Manufacturing standards for pipes commonly require a test of chemical composition and a series of mechanical strength tests for each heat of pipe. A heat of pipe is all forged from the same cast ingot, and therefore had the same chemical composition. Mechanical tests may be associated to a lot of pipe, which would be all from the same heat and have been through the same heat treatment processes. The manufacturer performs these tests and reports the composition in a mill traceability report and the mechanical tests in a material test report, both of which are referred to by the acronym MTR. Material with these associated test reports is called traceable. For critical applications, third party verification of these tests may be required; in this case an independent lab will produce a certified material test report(CMTR), and the material will be called certified.
Some widely used pipe standards or piping classes are:
The API range – now ISO 3183. E.g.: API 5L Grade B – now ISO L245 where the number indicates yield strength in MPa
ASME SA106 Grade B (Seamless carbon steel pipe for high temperature service)
ASTM A312 (Seamless and welded austenitic stainless steel pipe)
ASTM C76 (Concrete Pipe)
ASTM D3033/3034 (PVC Pipe)
ASTM D2239 (Polyethylene Pipe)
ISO 14692 (Petroleum and natural gas industries. Glass-reinforced plastics (GRP) piping. Qualification and manufacture)
ASTM A36 (Carbon steel pipe for structural or low pressure use)
ASTM A795 (Steel pipe specifically for fire sprinkler systems)
API 5L was changed in the second half of 2008 to edition 44 from edition 43 to make it identical to ISO 3183. It is important to note that the change has created the requirement that sour service, ERW pipe, pass a hydrogen induced cracking (HIC) test per NACE TM0284 in order to be used for sour service.
ACPA [American Concrete Pipe Association]
AWWA [American Water Works Association]
AWWA M45
Installation
Pipe installation is often more expensive than the material and a variety of specialized tools, techniques, and parts have been developed to assist this. Pipe is usually delivered to a customer or jobsite as either "sticks" or lengths of pipe (typically , called single random length) or they are prefabricated with elbows, tees and valves into a prefabricated pipe spool [A pipe spool is a piece of pre-assembled pipe and fittings, usually prepared in a shop so that installation on the construction site can be more efficient.]. Typically, pipe smaller than are not pre-fabricated. The pipe spools are usually tagged with a bar code and the ends are capped (plastic) for protection. The pipe and pipe spools are delivered to a warehouse on a large commercial/industrial job and they may be held indoors or in a gridded laydown yard. The pipe or pipe spool is retrieved, staged, rigged, and then lifted into place. On large process jobs the lift is made using cranes and hoist and other material lifts. They are typically temporarily supported in the steel structure using beam clamps, straps, and small hoists until the pipe supports are attached or otherwise secured.
An example of a tool used for installation for a small plumbing pipe (threaded ends) is the pipe wrench. Small pipe is typically not heavy and can be lifted into place by the installation craft laborer. However, during a plant outage or shutdown, the small (small bore) pipe may also be pre-fabricated to expedite installation during the outage. After the pipe is installed it will be tested for leaks. Before testing it may need to be cleaned by blowing air or steam or flushing with a liquid.
Pipe supports
Pipes are usually either supported from below or hung from above (but may also be supported from the side), using devices called pipe supports. Supports may be as simple as a pipe "shoe" which is akin to a half of an I-beam welded to the bottom of the pipe; they may be "hung" using a clevis, or with trapeze type of devices called pipe hangers. Pipe supports of any kind may incorporate springs, snubbers, dampers, or combinations of these devices to compensate for thermal expansion, or to provide vibration isolation, shock control, or reduced vibration excitation of the pipe due to earthquake motion. Some dampers are simply fluid dashpots, but other dampers may be active hydraulic devices that have sophisticated systems that act to dampen peak displacements due to externally imposed vibrations or mechanical shocks. The undesired motions may be process derived (such as in a fluidized bed reactor) or from a natural phenomenon such as an earthquake (design basis event or DBE).
Pipe hanger assembles are usually attached with pipe clamps. Possible exposure to high temperatures and heavy loads should be included when specifying which clamps are needed.
Joining
Pipes are commonly joined by welding, using threaded pipe and fittings; sealing the connection with a pipe thread compound, Polytetrafluoroethylene (PTFE) Thread seal tape, oakum, or PTFE string, or by using a mechanical coupling. Process piping is usually joined by welding using a TIG or MIG process. The most common process pipe joint is the butt weld. The ends of pipe to be welded must have a certain weld preparation called an End Weld Prep (EWP) which is typically at an angle of 37.5 degrees to accommodate the filler weld metal. The most common pipe thread in North America is the National Pipe Thread (NPT) or the Dryseal (NPTF) version. Other pipe threads include the British Standard Pipe Thread (BSPT), the garden hose thread (GHT), and the fire hose coupling (NST).
Copper pipes are typically joined by soldering, brazing, compression fittings, flaring, or crimping.
Plastic pipes may be joined by solvent welding, heat fusion, or elastomeric sealing.
If frequent disconnection will be required, gasketed pipe flanges or union fittings provide better reliability than threads. Some thin-walled pipes of ductile material, such as the smaller copper or flexible plastic water pipes found in homes for ice makers and humidifiers, for example, may be joined with compression fittings.
typically uses a "push-on" gasket style of pipe that compresses a gasket into a space formed between the two adjoining pieces. Push-on joints are available on most types of pipe. A pipe joint lubricant must be used in the assembly of the pipe. Under buried conditions, gasket-joint pipes allow for lateral movement due to soil shifting as well as expansion/contraction due to temperature differentials. Plastic MDPE and HDPE gas and water pipes are also often joined with Electrofusion fittings.
Large above ground pipe typically uses a flanged joint, which is generally available in ductile iron pipe and some others. It is a gasket style where the flanges of the adjoining pipes are bolted together, compressing the gasket into a space between the pipe.
Mechanical grooved couplings or Victaulic joints are also frequently used for frequent disassembly and assembly. Developed in the 1920s, these mechanical grooved couplings can operate up to working pressures and available in materials to match the pipe grade. Another type of mechanical coupling is a flareless tube fitting (Major brands include Swagelok, Ham-Let, Parker); this type of compression fitting is typically used on small tubing under in diameter.
When pipes join in chambers where other components are needed for the management of the network (such as valves or gauges), dismantling joints are generally used, in order to make mounting/dismounting easier.
Fittings and valves
Fittings are also used to split or join a number of pipes together, and for other purposes. A broad variety of standardized pipe fittings are available; they are generally broken down into either a tee, an elbow, a branch, a reducer/enlarger, or a wye. Valves control fluid flow and regulate pressure. The piping and plumbing fittings and valves articles discuss them further.
Cleaning
The inside of pipes can be cleaned with a tube cleaning process, if they are contaminated with debris or fouling. This depends on the process that the pipe will be used for and the cleanliness needed for the process. In some cases the pipes are cleaned using a displacement device formally known as a Pipeline Inspection Gauge or "pig"; alternately the pipes or tubes may be chemically flushed using specialized solutions that are pumped through. In some cases, where care has been taken in the manufacture, storage, and installation of pipe and tubing, the lines are blown clean with compressed air or nitrogen.
Other uses
Pipe is widely used in the fabrication of handrails, guardrails, and railings.
Applications
Steel pipe
Steel pipe (or black iron pipe) was once the most popular choice for supply of water and flammable gases. Steel pipe is still used in many homes and businesses to convey natural gas or propane fuel, and is a popular choice in fire sprinkler systems due to its high heat resistance. In commercial buildings, steel pipe is used to convey heating or cooling water to heat exchangers, air handlers, variable air volume (VAV) devices, or other HVAC equipment.
Steel pipe is sometimes joined using threaded connections, where tapered threads (see National Pipe Thread) are cut into the end of the tubing segment, sealant is applied in the form of thread sealing compound or thread seal tape (also known as PTFE or Teflon tape), and it is then threaded into a corresponding threaded fitting using two pipe wrenches. Beyond domestic or light commercial settings, steel pipe is often joined by welding, or by use of mechanical couplings made by companies such as Victaulic or Anvil International (formerly Grinnell) that hold the pipe joint together via a groove pressed or cut (a rarely used older practice), into the ends of the pipes.
Other variations of steel pipe include various stainless steel and chrome alloys. In high-pressure situations these are usually joined by TIG welding.
In Canada, with respect to natural gas (NG) and propane (LP gas), black iron pipe (BIP) is commonly used to connect an appliance to the supply. It must however be marked (either painted yellow or yellow banding attached at certain intervals) and certain restrictions apply to which nominal pipe size (NPS) can be put through walls and buildings. With propane in particular, BIP can be run from an exterior tank (or cylinder) provided it is well protected from the weather, and an anode-type of protection from corrosion is in place when the pipe is to be installed underground.
Copper pipe
Copper tubing is most often used for supply of hot and cold water, and as refrigerant line in HVAC systems. There are two basic types of copper tubing, soft copper and rigid copper. Copper tubing is joined using flare connection, compression connection, or solder. Copper offers a high level of resistance to corrosion, but is becoming very costly.
Soft copper
Soft (or ductile) copper tubing can be bent easily to travel around obstacles in the path of the tubing. While the work hardening of the drawing process used to size the tubing makes the copper hard/rigid, it is carefully annealed to make it soft again; it is therefore more expensive to produce than non-annealed, rigid copper tubing. It can be joined by any of the three methods used for rigid copper, and it is the only type of copper tubing suitable for flare connections. Soft copper is the most popular choice for refrigerant lines in split-system air conditioners and heat pumps.
Flare connections
Flare connections require that the end of a tubing section be spread outward in a bell shape using a flare tool. A flare nut then compresses this bell-shaped end onto a male fitting. Flare connections are a labor-intensive method of making connections, but are quite reliable over the course of many years.
Rigid copper
Rigid copper is a popular choice for water lines. It is joined using a sweat, compression or crimped/pressed connection. Rigid copper, rigid due to the work hardening of the drawing process, cannot be bent and must use elbow fittings to go around corners or around obstacles. If heated and allowed to slowly cool, called annealing, then rigid copper will become soft and can be bent/formed without cracking.
Soldered connections
Solder fittings are smooth, and easily slip onto the end of a tubing section. Both the male and female ends of the pipe or pipe connectors are cleaned thoroughly then coated with flux to make sure there is no surface oxide and to ensure that the solder will bond properly with the base metal. The joint is then heated using a torch, and solder is melted into the connection. When the solder cools, it forms a very strong bond which can last for decades. Solder-connected rigid copper is the most popular choice for water supply lines in modern buildings. In situations where many connections must be made at once (such as plumbing of a new building), solder offers much quicker and much less expensive joinery than compression or flare fittings. The term sweating is sometimes used to describe the process of soldering pipes.
Compression connections
Compression fittings use a soft metal or thermoplastic ring (the compression ring or "ferrule") which is squeezed onto the pipe and into the fitting by a compression nut. The soft metal conforms to the surface of the tubing and the fitting, and creates a seal. Compression connections do not typically have the long life that sweat connections offer, but are advantageous in many cases because they are easy to make using basic tools. A disadvantage in compression connections is that they take longer to make than sweat, and sometimes require retightening over time to stop leaks.
Crimped or pressed connections
Crimped or pressed connections use special copper fittings which are permanently attached to rigid copper tubing with a powered crimper. The special fittings, manufactured with sealant already inside, slide over the tubing to be connected. Thousands of pounds-force per square inch of pressure are used to deform the fitting and compress the sealant against the inner copper tubing, creating a watertight seal. Advantages of this method are:
A correctly crimped connection should last as long as the tubing.
It takes less time to complete than other methods.
It is cleaner in both appearance and the materials used to make the connection.
No open flame is used during the connection process.
Disadvantages are:
The fittings used are harder to find and cost significantly more than sweat type fittings.
The fittings are not re-usable. If a design change is required or if a joint is found to be defective or improperly crimped, the already installed fittings must be cut out and discarded. In addition, the cutting required to remove the fitting often will leave insufficient tubing to install the new fitting, So couplers and additional tubing will need to be installed on either side of the replacement fitting. Whereas with a soldered fitting, a defective joint can just be re-soldered, or heated and turned if a minor change is required, or heated and removed without requiring any of the tubing to be cut away. This also allows more expensive fittings like valves to be re-used if they are otherwise in good to new condition, something not possible if the fitting is crimped on.
The cost of the tooling is very expensive. , a basic toolkit required to sweat solder all the copper pipes of a typical single family residence, including fuel and solder, can be purchased for approximately $200. By contrast, the minimum cost of a basic powered crimping tool starts at around $1800, and can be as high as $4000 for the better brands with a complete set of crimping dies.
Aluminium pipe
Aluminium is sometimes used due to its low cost, resistance to corrosion and solvents, and its ductility. Aluminium tube is more desirable than steel for the conveyance of flammable solvents, since it cannot create sparks when manipulated. Aluminium tubing can be connected by flare or compression fittings, or it can be welded by the TIG or heliarc processes.
Glass pipe
Tempered glass pipes are used for specialized applications, such as corrosive liquids, medical or laboratory wastes, or pharmaceutical manufacturing. Connections are generally made using specialized gasket or O-ring fittings.
Plastic pipe
Plastic pipes used in manufacturing.
Plastic pipe fittings include PVC pipe fittings, PP / PPH pipe fitting mould, PE pipe and ABS pipe fitting.
| Technology | Components_2 | null |
3705573 | https://en.wikipedia.org/wiki/Gorgonopsia | Gorgonopsia | Gorgonopsia (from the Greek Gorgon, a mythological beast, and 'aspect') is an extinct clade of sabre-toothed therapsids from the Middle to the Upper Permian, roughly between 270 and 252 million years ago. They are characterised by a long and narrow skull, as well as elongated upper and sometimes lower canine teeth and incisors which were likely used as slashing and stabbing weapons. Postcanine teeth are generally reduced or absent. For hunting large prey, they possibly used a bite-and-retreat tactic, ambushing and taking a debilitating bite out of the target, and following it at a safe distance before its injuries exhausted it, whereupon the gorgonopsian would grapple the animal and deliver a killing bite. They would have had an exorbitant gape, possibly in excess of 90°, without having to unhinge the jaw.
They markedly increased in size as time went on, growing from small skull lengths of in the Middle Permian to bear-like proportions of up to in the Upper Permian. The latest gorgonopsians, Rubidgeinae, were the most robust of the group and could produce especially powerful bites. Gorgonopsians are thought to have been completely terrestrial and could walk with a semi-erect gait, with a similar terrestrial locomotory range as modern crocodilians. They may have been more agile than their prey items, but were probably inertial homeotherms rather than endotherms unlike contemporary therocephalians and cynodonts, and thus were probably comparatively less active. Though gorgonopsians were able to maintain a rather high body temperature, it is unclear if they would have also had sweat glands or fur (and by extension whiskers and related structures). Their brains were reminiscent of modern reptilian brains, rather than those of living mammals. Most species may have been predominantly diurnal (active during the day) though some could have been crepuscular (active at dawn or dusk) or nocturnal (active at night). They are thought to have had binocular vision, a parietal eye (which detects sunlight and maintains circadian rhythm), a keen sense of smell, a functional vomeronasal organ ("Jacobson's organ"), and possibly a rudimentary eardrum.
The major therapsid groups had all evolved by 275 million years ago from a "pelycosaur" ancestor (a poorly defined group including all synapsids which are not therapsids). The therapsid takeover from pelycosaurs took place by the Middle Permian as the world progressively became drier. Gorgonopsians rose to become apex predators of their environments following the Capitanian mass extinction event which killed off the dinocephalians and some large therocephalians after the Middle Permian. Despite the existence of a single continent during the Permian, Pangaea, gorgonopsians have only been found in the Karoo Supergroup (primarily in South Africa, but also in Tanzania, Zambia, and Malawi), the Moradi Formation of Niger, western Russia, and in the Turpan Basin of Xinjiang, China, with probable remains known from the Kundaram Formation in the Pranhita–Godavari Basin of India. These places were semi-arid areas with highly seasonal rainfall. The oldest gorgonopsian specimen, a middle Permian gorgonopsian, was described from the Port des Canonge Formation on the island of Mallorca, western Mediterranean. Gorgonopsian genera are all very similar in appearance, and consequently many species have been named based on flimsy and likely age-related differences since their discovery in the late 19th century, and the group has been subject to several taxonomic revisions.
Most gorgonopsians became extinct during a phase of the Permian–Triassic extinction event taking place at the very end of the Permian, in which major volcanic activity (which would produce the Siberian Traps) and resultant massive spike in greenhouse gases caused rapid aridification due to temperature spike, acid rain, frequent wildfires, and potential breakdown of the ozone layer. However, some smaller taxa like Cyonosaurus may have survived up to the Early Triassic. The large predatory niches would be taken over by the archosaurs (namely crocodilians and dinosaurs) in the Mesozoic.
Description
Earlier gorgonopsids in the Middle Permian were quite small, with skull lengths of , whereas some later genera attained massive, bear-like sizes with the largest being Inostrancevia up to in length and in body mass. Nonetheless, small gorgonopsians remained abundant until extinction (though small species may actually represent juvenile specimens of other taxa).
Like other Permian therapsids, gorgonopsians had developed several mammalian characteristics. These might have included a parasagittal gait (the limbs were vertically oriented and moved parallel to the spine) as opposed to the sprawling gait of amphibians and earlier synapsids. This gait change in therapsids was possibly related to the reduction in tail size and phalangeal formula (the number of bones per digit, which for gorgonopsians was 2.3.4.5.3 like reptiles). Other developments included fibrous lamellar cortical bone and deeply-set teeth. Like reptiles, gorgonopsians lack a secondary palate separating the mouth from the nasal cavity, prohibiting chewing.
Skull
Anatomy varies incredibly little between gorgonopsians. Many species are distinguished by vague proportional differences, and consequently smaller species may actually represent juveniles of larger taxa. Notably, the vomer at the tip of the snout varies among species in terms of the degree of its expansion, as well as the positions, degree of splay, and shape of the 3 ridges. They typically feature a long and narrow skull. Juvenile Rubidgea appear to have had snouts wider than long. Unlike eutheriodonts, the occipital bone (at the back of the skull) is rectangular (wider than tall) and concave, as opposed to triangular.
The gorgonopsian brain, like other non-mammaliaform therapsids, lacks an expansion of the neocortex, has a relatively large hindbrain compared to the forebrain, a large epyphysial nerve (found in creatures with a parietal eye on the top of the head), an enlarged pituitary gland, and an overall elongated shape; all-in-all resembling a reptilian brain. The braincase was also rather reptilian, and is also comparatively smaller and not as thick as those of mammals. The flocculus, a lobe of the cerebellum, is proportionally large, and is related to the vestibulo-ocular reflex (which stabilises gaze while moving the head). Judging by the orientation of the semi-circular canals in the ear (which have to be oriented parallel to the ground), the head of the gorgonopsian specimen GPIT/RE/7124 would have tilted forward by about 41°, increasing the overlap between the visual fields of the two eyes and improving binocular vision—useful to a predator. Unlike either reptiles or mammals, the semi-circular canals are flat, probably because they were wedged between the opisthotic (an inner ear bone) and supraoccipital bones.
Teeth
Like many mammals, gorgonopsians were heterodonts, with clearly defined incisors, canines, and postcanine teeth homologous with premolars and molars. They had five incisors in the upper jaw (for most, the first three were the same size as each other, and the last two were shorter) and four on the bottom.
In the majority of gorgonopsians, the incisors were large, and the upper canines were elongated into sabres, much like those of later sabre-toothed cats. Some gorgonopsians had exceptionally long upper canines, such as Inostrancevia, and some of them had a flange on the lower jaw to sheath the tip of the canine while the mouth was closed. Sabres are generally interpreted as having been used as stabbing or slashing weapons, which would have required an extremely wide gape. Both the upper and lower canines of Rubidgea were elongated, and the animal would have needed an even greater gape. The serration pattern of gorgonopsians was most similar to those of theropod dinosaurs than to other synapsids. The palate also features tuberosities and ridges which oftentimes have functional teeth, which may have been used to hold onto struggling prey, diverting these powerful forces away from the fragile canines. Similar ridges have been identified on the machairodont Homotherium. The postcanine teeth were reduced in both size and number; many rubidgeines (the latest gorgonopsians) did not have postcanines in the lower jaw, and Clelandina lacked them entirely.
Gorgonopsians were polyphyodonts, and teeth grew continuously throughout an individual's life. Like some therapsids, while there was one functional canine, another canine was growing to replace it when it inevitably broke off. The left and right sides of the jaws did not have to be synchronous, so, for example, the first canine on the left side could be functional while the first canine on the right side was still growing. Such a method might have been in play so as always to have a set of functional canines, as having a single or no canines would have severely impeded hunting, and growing such large teeth took a long time. On the other hand, because the functional canine is typically found in the foremost tooth socket (instead of equal occurrence in either socket), it is possible that canine replacement occurred a finite number of times, and the animal would eventually be left with a single, permanent set of functional canines in these sockets. In 1984, British palaeontologists Doris and Kenneth Kermack suggested that the canines grew to match the size of the skull, and continually broke off until the animal stopped growing, and that gorgonopsians featured an early version of finite tooth replacement exhibited in many mammals. The tooth replacement patterns of the other teeth are unclear. The postcanine teeth were replaced more slowly than the other teeth, likely due to their lack of functional significance.
Postcranium
The seven cervical vertebrae (in the neck) are all the same size as each other except for the last one, which is shorter and lower; there is one atlas and one axis. Like sabre-toothed cats, the neck is long with well-developed muscles, which would have been especially useful when the canines were sunk into an animal. Like other early synapsids, gorgonopsians have a single occipital condyle, and the articulation (the joints) of the cervical vertebrae is overall reptilian, permitting side-to-side movement of the head but restricting up-and-down motion. The last cervical is shaped more like the dorsal vertebrae.
The dorsals are spool-shaped and all appear about the same as each other. The spinous processes jut out steeply from the centra, and feature sharp keels on the front and back sides. Unlike eutheriodonts, gorgonopsians do not have distinguished lumbar vertebrae. Nonetheless, the dorsals equating to that series are similar to the lumbars of sabre-toothed cats with steeply oriented zygopophases, useful in stabilising the lower back especially when pinning down struggling prey.
There are three sacral vertebrae, and the series attached to the pelvis by the first vertebra. The pelvis is reptilian, with separated ilium, ischium, and pubis. The femur is slightly s-shaped, and is short but longer and slenderer than the humerus. For most, the tibia and fibula strongly curve into each other, and the tibia is more robust than the fibula. The joint between the ankle and the heel bones may have been somewhat mobile. The fifth digit for both the hands and feet was not attached to the carpus/tarsus, and instead connected directly to the ulna/heel bone.
Taxonomy
Fossil bearing sites
In 1876, the first gorgonopsian remains were identified in the Beaufort Group of the Karoo Supergroup of South Africa, by the biologist and paleontologist Richard Owen. He classified the fossils as Gorgonops torvus, combining the Greek Gorgon, a mythological beast, with the word (), meaning 'aspect'. In Africa, gorgonopsians have also been found in Karoo outcroppings in the Ruhuhu Valley of Tanzania, the Upper Luangwa Valley of Zambia, and Chiweta, Malawi.
Gorgonopsians were first identified in Russia in the 1890s at the Sokolki locality on the Northern Dvina in Siberia under the supervision of Russian palaeontologist Vladimir Prokhorovich Amalitskii. In a posthumous publication, it was described as Inostrancevia alexandri, and it is one of the best known and largest gorgonopsians. Since then, only a few more Russian taxa have been described: Pravoslavlevia, Sauroctonus, Viatkogorgon, Suchogorgon, and Nochnitsa. One of the Russian genera named, Leogorgon, is now considered as a nomen dubium.
Gorgonopsians are conspicuously rare beyond these two areas. In 1979, Chinese palaeontologist Yang Zhongjian described a Chinese gorgonopsian "Wangwusaurus tayuensis" based on teeth from the Late Permian Jiyuan Formation, but in 1981, palaeontologists Denise Sigogneau-Russell and Ai-Lin Sun found the assigned material to be a random assemblage of which only two have even a remote similarity to Gorgonopsia. In 2011, an incomplete set of teeth were discovered in a locality within the Turpan Depression, Xinjiang, in the northwest of the country. The fossil material, although thin, is described in 2022 by paleontologists Jun Liu and Wan Yiang and confirms that it comes from a gorgonopsian dating from the Upper Permian that actually lived in present-day China. In 2003, Indian palaeontologists Sanghamitra Ray and Saswati Bandyopadhyay assigned some skull fragments from the Late Permian Kundaram Formation to a medium-sized gorgonopsian, though the gorgonopsian characteristics have also been documented in some therocephalians. In 2008, a large and probably rubidgeine upper jaw fragment and canine was identified at the Late Permian Moradi Formation in Niger (one of the few low-latitude Late Permian tetrapod-bearing formations), and is the first evidence of a low-latitude gorgonopsian. A second low-latitude gorgonopsian was described in 2024 from the Port des Canonge Formation of Mallorca in the western Mediterranean. Notably, it also likely represents the oldest gorgonopsian yet known in the fossil record, constrained to a minimum age before the middle Wordian, if not pre-Wordian. This suggests that the absence of gorgonopsians from low-latitudes reflects uneven sampling rather than a true restriction to high latitudes.
Classification
Upon discovery, Owen presumed that Gorgonops and several other taxa he described from the Karoo Supergroup were cold-blooded reptiles, despite bearing teeth resembling those of carnivorous mammals. He proposed classifying all of them under the newly coined order Theriodontia (which he placed in the class Reptilia). He decided to subdivide Theriodontia into families based on the anatomy of the nostrils (the bony narials)—"Mononarialia" for those with one opening in the skull for the nose as in mammals, "Binarialia" for those with two openings as in reptiles, and "Tectinarialia" for Gorgonops because its opening was overshadowed by a thick bone roof (tectus is Latin for "covered, roofed, decked"). In 1890, English naturalist Richard Lydekker made Gorgonops the type species of the family Gorgonopsidae. British palaeontologist Harry Seeley in 1894 believed Gorgonops lacked an opening in the temporal bone (temporal fenestra), which is a diagnostic feature of Theriodontia, and so elevated Gorgonopsidae to Gorgonopsia, distinct from Theriodontia. He classified all South African materials bearing both reptilian and mammalian traits into the order "Theriosuchia", and considered Gorgonopsia and Theriodontia suborders of it. American palaeontologist Henry Fairfield Osborn completely reworked the classification of Reptilia in 1903, and erected two major groups: Diapsida and Synapsida, and in 1905, South African palaeontologist Robert Broom created a third group, Therapsida, to house the "mammal-like reptiles", including Theriodontia. He also challenged Seeley's claim and relegated Gorgonops back to Theriodontia, but he placed it into his newly erected subgroup Therocephalia, dissolving Gorgonopsia. In 1913, especially in light of an almost complete G. torvus skull discovered by the Reverend John H. Whaits, Broom reinstated Gorgonopsia.
The number of South African genera rapidly grew in the 20th-century, headed principally by Broom, whose extensive work on the Karoo therapsids—from the beginning of his career in the country in 1897 to his death in 1951—led to his description of 57 gorgonopsian holotype specimens and 29 genera. Many of Broom's taxa would later be invalidated. Many other contemporary workers created wholly new species or genera based on single specimens. Consequently, Gorgonopsia has been the subject of much taxonomic turmoil, and is one of the most problematic synapsid groups. Because the skull anatomy differs very little across taxa, many are defined based on vague proportional differences, including even the well-known members. Nominal species are distinguished predominantly by traits which are known to be quite variable depending on the age of the individual, including eye orbit size, snout length, and number of postcanine teeth. Thus, it is possible that some taxa are synonymous with each other, and represent different stages of development.
Among the first attempts to organise the clade was carried out by British zoologist David Meredith Seares Watson and American palaeontologist Alfred Romer in 1956, who split it into twenty families, of which the members of three (Burnetiidae, Hipposauridae, and Phthinosuchidae) are not considered gorgonopsians anymore. In 1970 and again in 1989, predominantly considering African taxa, Sigogneau-Russell published a comprehensive monograph on Gorgonopsia (defining it as an infraorder), and recognised only two families: Watongiidae and Gorgonopidae. Watongia was moved to Varanopidae in 2004 by the Canadian paleontologists Robert R. Reisz and Michel Laurin. Sigogneau-Russell split Gorgonopidae into three subfamilies—Gorgonopsinae, Rubidgeinae, and Inostranceviinae—and reduced the number of genera to twenty-three. In 2002, Russian palaeontologist Mikhail Feodosievich Ivakhnenko, considering the Russian taxa, instead considered Gorgonopsia a suborder, and grouped it together with Dinocephalia into the order "Gorgodontia". He divided Gorgonopsia into the superfamilies "Gorgonopioidea" (families Gorgonopidae, Cyonosauridae, and Galesuchidae) and "Rubidgeoidea" (Rubidgeidae, Phtinosuchidae, and Inostranceviidae). In 2007, biologist Eva V. I. Gebauer, in her comprehensive review of Gorgonopsia (her PhD dissertation), rejected Ivakhnenko's model in favour of Sigogneau-Russell's, and further reduced the number of genera to fourteen in addition to the Russian genera: Aloposaurus, Cyonosaurus, Aelurosaurus, Sauroctonus, Scylacognathus, Eoarctops, Gorgonops, [["Dixeya" nasuta|"Dixeya" nasuta]] (under the informal nomen nudum "Njalila"), Lycaenops, Arctognathus, Aelurognathus, Sycosaurus, Clelandina, and Rubidgea. In general, Sigogneau-Russell's model is supported, but there is little consensus on which genera can be assigned to which subfamilies. In 2015, American palaeontologist Christian F. Kammerer and colleagues redescribed Eriphostoma (which was labelled as an indeterminate theriodont) as a gorgonopsian, and sunk Scylacognathus and the next year Eoarctops into it.
The first phylogeny (family tree) of the members of Gorgonopsia was published in 2016 by American palaeontologist Christian F. Kammerer, who specifically investigated Rubidgeinae, and re-described both the subfamily and the nine species he assigned to it (reducing the number from thirty-six species). Kammerer also resurrected Dinogorgon, Leontosaurus, Ruhuhucerberus, and Smilesaurus. Kammerer was unsure if Leontosaurus, Clelandina, Dinogorgon, and Rubidgea all represent the same taxon or not (for which Dinogorgon has priority), but he decided to classify all of them in the tribe Rubidgeini pending further examination. In 2018, Kammerer and Russian palaeontologist Vladimir Masyutin identified a new genus Nochnitsa as the basalmost known gorgonopsians, and found that all Russian taxa (except Viatkogorgon, which is in the outclade) form a completely separate clade from the African taxa. Also in 2018, palaeobiologist Eva-Maria Bendel, Kammerer, and colleagues resurrected Cynariops. In 2022, Kammerer and fellow palaeontologist Bruce S. Rubidge described Phorcys from South Africa.
Evolution
Synapsida has traditionally been split into the basal "Pelycosauria" and the derived Therapsida. The former comprises cold-blooded creatures with a sprawling gait and presumably lower metabolism which evolved in the Upper Carboniferous. Through the middle to late 20th-century, American palaeontologist Everett C. Olson investigated synapsid diversity in the Middle Permian San Angelo, Flowerpot, and Chickasha Formations in North America, and noted that pelycosaur diversity reduced from six to three in these formations, and that they coexisted with several fragmentary specimens which he interpreted as therapsids. He then suggested the adaptive shift from pelycosaur-grade to therapsid-grade took place during the Middle Permian (Olson's Extinction); however, the classification of those "therapsids" and the age of the formations have since been challenged. Thus, the exact timing of the therapsid takeover is unclear, but the six major therapsid groups (Biarmosuchia, Dinocephalia, Anomodontia, Gorgonopsia, Therocephalia, and Cynodontia) had evolved by 265 million years ago during the Wordian.
The oldest definitive gorgonopsian fossil worldwide comes from the Port des Canonge Formation of Mallorca in the western Mediterranean. The rocks of this formation are imprecisely dated, however magnetostratigraphy, palynology and ichnology constrain its age to the earliest Wordian at minimum. The stratigraphic position of the gorgonopsian itself and surrounding fossil track record may suggest it is significantly older, potentially Roadian of the early Guadalupian (Middle Permian) or even of the Kungurian or Artinskian stages of the Cisuralian (Early Permian). Such an age would not only make it the oldest gorgonopsian yet known, but the oldest therapsid altogether. Although of indeterminable species, the anatomy of the Mallorcan gorgonopsian suggests it is more derived than the earliest-diverging recognised gorgonopsian, Nochnitsa, and is consistent with a position in or just outside of the base of the "African" and "Russian" clades. This suggests that not only were gorgonopsians present in the earliest Wordian, if not the late Cisuralian, but that their diversification was well underway by this time and that the clade as a whole likely originated before this interval. Furthermore, it suggests that early gorgonopsian (and indeed therapsid) diversity was not restricted to high-latitudes of temperate Pangaea, as previously suggested from literal readings of the fossil record, but also included the low equatorial latitudes of the supercontinent (where Mallorca was located).
The oldest high-latitude gorgonopsians both come from the Karoo Sugergroup of South Africa. The older of the two is a partial snout belonging to a large undeterminable genus from the Eodicynodon Assemblage Zone of the Karoo Basin, roughly dating to the Wordian. Phorcys, known from the lowermost end of Karoo's Tapinocephalus Assemblage Zone (roughly dating a little later to the Wordian/Capitanian boundary), is the oldest identifiable gorgonopsian taxon. Although highly fragmentary, both Karoo taxa are estimated to have had skulls measuring approximately in length, significantly larger than the estimated ~ skull of the Mallorcan gorgonopsian and the slightly younger Eriphostoma. This suggests gorgonopsian body size was also diversifying early in their evolution.
The Permian progressively became dryer and dryer. In the Upper Carboniferous and Lower Permian, pelycosaurs seem to have clung to the everwet coal swamp habitats near the equator (fossils known within 10° of either side of the palaeoequator); beyond this to about 30° was an expansive desert which extended all the way to the coast, separating the swamps from the temperate regions. By the Middle Permian, the equatorial forests had switched to a seasonal wet/dry system, but the swamps were connected to the temperate zones via coastal passages along East Pangaea, allowing cross-continental migration from what is now South Africa to what is now Russia. Therapsids appear to have evolved in this seasonally humid/dry landscape, expanding even into the temperate zones. At this point, synapsids were the only large terrestrial animals of their environment; and pelycosaurs may not have been able to adapt to the aridification. At about the time of pelycosaur extinction, therapsids experienced a major adaptive radiation (all carnivores) continuing into the Upper Permian.
Throughout the Middle Permian, the often gigantic dinocephalians were the dominant animals of their ecosystems. They disappear from the fossil record during the Capitanian mass extinction event caused by volcanic activity which has formed the Chinese Emeishan Traps. The exact cause of their extinction is unclear, but they were replaced by gorgonopsians and dicynodonts (which began to greatly increase in size) and the smaller therocephalians. The rubidgeans were the most derived gorgonopsians, and consequently the most massive and heavily built.
Palaeobiology
Bite
Gorgonopsians were likely active predators. The rubidgeines have an especially robust skull among gorgonopsians, comparable to those of enormous macropredators which use their skulls as their primary weapon, such as mosasaurs or some theropod dinosaurs. Less robust gorgonopsians with longer canines and much weaker bite, such as Smilesaurus or Inostrancevia, instead probably used their canines for slashing, much more similar to sabre-toothed cats. The postcanines of Clelandina were replaced by a smooth ridge unlike dicynodonts which have a blade-like keratinous ridge, and it may have predominantly gone after prey it could swallow whole. Gorgonopsian taxa did coexist with each other—as many as seven at one time—and the fact that some rubidgeines possess postcanines while some other contemporary ones do not suggests that they practiced niche partitioning and pursued different prey items.
The elongated canines have generally been thought to have been instrumental in their hunting tactics. The gorgonopsian jaw hinge was double jointed and made up of somewhat mobile and rotatable bones, which would have allowed them to open their mouths incredibly wide—perhaps in excess of 90°—without having to unhinge the jaw. It has alternatively been suggested (first in 2002 by biologists Blaire Van Valkenburgh and Tyson Secco, though in reference to cats) that sabres evolved primarily due to sexual selection as a form of mating display. This is exhibited in some modern deer species, but is difficult to test given the lack of living sabre-toothed synapsid predators. In sabre-toothed cats, long-sabred ("dirk-toothed") taxa are thought to have been pursuit hunters, whereas short-toothed ("scimitar-toothed") taxa are thought to have been ambush predators. Among the dirk-toothed cats, these predators are suggested to have killed with a well-placed slash to the throat after grappling prey, but gorgonopsians may have been less precise with bite placement, armed with reptilian jaws and tooth arrangements. Instead, gorgonopsians possibly used a bite-and-retreat tactic: the predator would ambush its quarry and take a sizable and debilitating bite out of it, and then follow as the prey tried to escape before succumbing to its injury, whereupon the gorgonopsian would deliver a killing bite. Because the postcanines are reduced or entirely absent, meat would have been forcibly torn away from the carcass and swallowed whole. This "puncture–pull" strategy is also hypothesised to have been used by theropod dinosaurs.
Gorgonopsians, along with other early carnivores as well as crocodiles, predominantly relied on "Kinetic-Inertial system" (KI) of biting down onto prey, in which the pterygoid and temporalis muscles rapidly clamped the jaws shut, using momentum and the kinetic energy of the jaws and teeth to grapple the victim. Mammalian carnivores, including sabre-toothed cats, instead rely mainly on the "Static-Pressure system" (SP) where the temporalis and masseter muscles produce a strong bite force to kill prey. The temporalis and masseter had only separated in mammals, and gorgonopsians instead had a muscle stretching from the underside of the skull roof, back to the squamosal bone (at the back of the skull), and across the cheekbones. The part anchored by the cheeks stabilised the jawbone and allowed it to move side-to-side while closing. This may have been very important in biting, as the cheekbones get stronger in tandem with the canines getting longer.
Smaller gorgonopsians, such as Cyonosaurus (which may actually represent a juvenile of a different species), had gracile skulls and sabres, and may have acted much like jackals and foxes. Bigger gorgonopsians, such as Gorgonops, had long robust snouts with strongly flared cheeks, which would have supported strong pterygoids and a powerful KI bite. The medium-size Arctognathus had a box-like skull and resultantly powerful snout, which would have allowed strong bending and torsion movements, and a combination of both KI and SP bite elements. Even bigger gorgonopsians, such as Arctops, had a shorter and more convex snout like the earlier sphenecodont Dimetrodon, and would have been able to rapidly clamp the jaws shut from a wide gape (which would have been necessary given the long canines). The even larger Rubidgeinae had extremely powerful, heavily built, buttressed skulls, with wide snouts, strongly flanged cheeks, and exceedingly long teeth; the sabres of Rubidgea atrox are longer than the teeth of Tyrannosaurus.
Unlike mammalian carnivores, gorgonopsians (and therocephalians) had reduced or completely lacked postcanines, and the jaw likely could not exert shearing pressure necessary for crushing bones open to access the bone marrow. It has largely been unclear if bone marrow had even evolved yet in Permian synapsids (fish and many amphibians lack this in present day), but in 2021 it was shown that the Early Permian amphibians Seymouria and Discosauriscus likely had haematopoietic (red blood cell-producing) bone marrow in their limbs.
Locomotion
Gorgonopsians are considered to have been strictly terrestrial. They are thought to have been able to move with an erect gait similar to that used by crocodilians, the limbs positioned almost vertically as opposed to horizontally as in the sprawling gait of lizards. The glenoid cavity on the shoulder blade is strongly angled tailwards, so the limbs had limited forward movement, and they may have had a short stride length. Lizards often move their spines side to side to increase stride length, but the more vertically orientated facet joints connecting the vertebrae in gorgonopsians would have made the spine more rigid and stable, encumbering such movement.
The gorgonopsian shoulder joint has a highly unusual configuration. The humeral head which connects to the shoulder is longer than the glenoid, so it could not fit into the cavity. Consequently, they may have been attached with a large mass of cartilage, with the humerus performing a rolling movement over the glenoid. This could theoretically make the angle between the humerus and the glenoid anywhere from 80 to 145° when facing the animal. If the angle was on the lower end, this would have been a rather firm joint, allowing the deltoids to exert great force through the forelimb, such as when pinning down struggling prey, or holding down a carcass while ripping off flesh. If the humerus was positioned at a higher angle, this could have permitted enhanced extension forwards and backwards (along the long axis) and thus greater stride length, useful in an attack or short chases. The shoulder blade expands off to the sides of the animal (protrudes laterally), also providing a large attachment for the deltoids. All the scapulohumeral muscles had strongly developed attachments, particularly the deltoids. When extending the forelimbs, the deltoids may have raised the front side (anterior margin) of the humerus, and coracobrachialis muscle lowered the back side (posterior margin). When retracting the forelimb, the pectoralis muscle may have pushed the anterior margin down, and the subscapularis muscle pulled the posterior margin up.
The pelvis joint has the usual ball-and-socket joint configuration. The somewhat flattened femoral head could theoretically have fit into the hip socket at a wide range of angles. In 1982, palaeontologist Tom S. Kemp suggested that early theriodonts, including gorgonopsians, could place the femur at both a horizontal angle in a sprawling gait, as well as a more vertical angle in an erect gait. He compared the locomotory habits of these creatures to those of crocodilians, which utilise a sprawling gait over short distances, but switch to an erect one while running or moving over longer distances. Though the hip of the specimen GPIT/RE/7113 seems to be anatomically intermediate between Dimetrodon and mammals—with the ilium expanded more in the headwards direction than the tailwards, and the pubis somewhat reduced—the puboischiofemoralis muscle (a large muscle carried only by reptiles which runs from the pelvis to the femur) extensively attached to the underside of the pubis and ischium, which would have allowed it to produce a strong adducting force (drawing the legs closer to the body), useful in a sprawling gait. It is also conceivable that gorgonopsians primarily engaged this muscle while grappling struggling prey. The shins are relatively short compared to the femur, which suggests gorgonopsians were not well adapted for running long distances.
In regard to how the feet were placed on the ground, gorgonopsians are the only early therapsids which present ectataxony (the last digit bears the most weight), homopody (footprints and handprints look the same), and semi-plantigrady (to some degree, the feet were placed flat on the ground). These adaptations may have made gorgonopsians swifter and more agile than their prey. Gorgonopsians had rather nimble digits, indicative of grasping capability for both the hands and feet, possibly for grappling struggling prey to prevent excessive load bearing on, and consequential fracturing or breaking of, the canines while they were sunk into the victim.
Senses
Unlike eutheriodonts, but like some ectothermic creatures today, all gorgonopsians possessed a pineal eye on the top of the head, which is used to detect daylight (and thus, the optimal temperature to be active). It is possible that other theriodonts lost this due to the evolution of either endothermy, intrinsically photosensitive retinal ganglion cells in the eyes—in tandem with the loss of colour vision and a shift to nocturnal life—or both. Nocturnal behaviour has long been assumed to have originated in mammals (nocturnal bottleneck), but the large orbit size and presence of sclerotic rings in many early synapsids, stretching as far back as the Carboniferous, would suggest that the ability to venture out in low-light conditions evolved much earlier. Based on these aspects, the specimen SAM-PK-K10034 may have had mesopic vision, and Cyonosaurus scotopic or photopic vision. The diameters of the sclerotic rings for the small Viatkogorgon are proportionally large, with an inner diameter of and outer diameter of , compared to a diameter of for the orbit itself, which suggests it made predominantly nocturnal excursions. Among gorgonopsians, the rubidgeine Clelandina has unusually small sclerotic rings, indicating it had photopic vision and was strictly diurnal; Kammerer suggested that niche partitioning among rubidgeines (as there have been as many as seven different taxa coexisting in an area), in part, took the form of different species being active at different times of the day, but the sclerotic rings of only Clelandina among this subfamily have been identified, making this hypothesis highly speculative.
Gorgonopsians have a rather short nasal cavity, like pelycosaurs, but it features abundant longitudinal ridges behind the internal nostrils (which connect the nasal cavity to the throat); because respired air would not have passed through them, these are typically interpreted as having been olfactory turbinates, and would have given gorgonopsians a rather highly developed sense of smell. Gorgonopsians possessed a vomeronasal organ ("Jacobson's organ")—a part of the accessory olfactory system—which would have been placed at the base of the nasal septum; unlike dicynodonts and therocephalians, there seems to have been a canal connecting the organ with the mouth, indicating it was functional in gorgonopsians.
Early theriodonts (including gorgonopsians) may have possessed an eardrum, unlike earlier pelycosaurs, indicated by the reduction of the connection between the quadrate bone (at the jaw hinge) and the pterygoid bone (at the palate), allowing the quadrate to independently vibrate to a degree. This may have allowed the detection of air-borne sounds with a low amplitude of less than , but the eardrum would have been supported by cartilage or ligaments instead of bone. If correct, then the postdentary bones (which in early mammals form the middle ear bones) would have needed to become detached from the dentary (jawbone); the gorgonopsian fossil record seems to indicate the postdentary-dentary connection was reduced. Though, given the specialisations required for biting, the condition of an isolating quadrate in gorgonopsians could alternatively be explained as streptostyly (rotatable quadrate) in order to widen the gape rather than facilitate hearing.
Thermoregulation
A major anatomical shift occurred between earlier pelycosaurs and therapsids, which is postulated to have been related to an increasing metabolism and the origins of homeothermy (maintenance of a high body temperature). The evolution of a secondary palate, and the separation of the mouth from the nasal cavity, may have increased ventilation efficiency associated with high levels of aerobic activity; gorgonopsians did not have a bony secondary palate, but possibly had one of soft tissue. Nonetheless, the secondary palate could have instead aided in eating large quantities of food at once rather than in ventilation. The reorganisation of the skeleton (from a sprawling to a parasagittal gait) has been postulated to be indicative of the presence of a diaphragm, and thus also enhanced ventilation for aerobic activity; but it could have instead been to increase acceleration or agility, which does not necessarily equate to intense aerobic activity, much like in crocodiles. Fibrous lamellar cortical bone, which all early therapsids had, would indicate an increased growth rate, but this may not be linked to metabolic rate.
Modern large reptiles naturally give off body heat at a slower rate than smaller ones, and are considered "inertial homeotherms", but they maintain a low body temperature of . If therapsids required a higher body temperature of , they would either have needed to have been endotherms (generating their own body heat) or have had greater control over heat loss (that is, better homeothermy). The parasagittal gait may have aided the latter, as it would have kept most of the body off the ground as well as allowed blood to stay in the abdomen instead of having to circulate through the appendages, both of which would reduce heat transfer to the ground and stabilise core temperature. The reduced tail would have also reduced the total surface area of the animal, further minimising heat loss. Among therapsids, only eutheriodonts (not gorgonopsians) have respiratory nasal turbinates, which help retain moisture while breathing in large quantities of air, and its evolution is typically associated with the beginning of "mammalian" oxygen consumption rates and the origins of endothermy.
If gorgonopsians were inertial homeotherms, it is not impossible that they had hair. The snout is typically riddled with foramina (small holes which confer with blood vessels), which could potentially point to the existence of loose skin (as opposed to scales), hair, various skin glands (such as sweat glands), and whiskers; however, some reptiles present a similar patterning of foramina, which are instead related to dental development rather than skin.
Palaeopathology
The anterolateral aspect of the left radius (a forearm bone) of the gorgonopsian specimen NHCC LB396 presents a circular bony lesion, featuring irregular-to-radial spikes made of cortical bone surrounded by a thin layer of subperiosteal bone, which grew rapidly over a single growing season. This is consistent with periostitis most likely stemming from subperiosteal haematoma. This specific condition as well as the fast growth rate are more reminiscent of mammals and dinosaurs than crocodilians or monitor lizards. Among early synapsids, the only other pathology noted is osteomyelitis in several pelycosaur groups.
The labial (lip/cheek) side of the tooth root of a functional canine of RB382 presents as many as 8 lesions, clustering along the midline of the tooth, which resemble miniature teeth with a pulp, dentine, and a thin enamel coating. They are roughly circular—with diameters varying from —though they become less circular at around the middle point of the root until passing the cervix of the tooth. This is roughly consistent with the human ailment odontoma, the most frequent type of odontogenic tumour, which previously only extended a few million years back in the fossil record. At 255 million years old, RB382 presents the oldest-known case of odontoma.
The adult snout SAM-PK-11490 from an indeterminate Middle Permian gorgonopsian species has an imbedded tooth from an unidentifiable animal. The bone developed a callus around the tooth, indicating it healed and the individual survived the attack. It either came from a predator—namely a biarmosuchian, a therocephalian, or another gorgonopsian—or intraspecific face biting as is commonly exhibited in social predators—such as big cats or monitor lizards, and it has been suggested for several extinct lineages such as theropods, aquatic reptiles, and saber-toothed cats. Social biting is intended to assert dominance or facilitate breeding, and, if correct, suggests at least some Middle Permian gorgonopsians were social carnivores. The tooth was initially overlooked so it is unclear how common this pathology actually is.
Palaeoecology
Paleoenvironment
Following the extinction of the dinocephalians and (in South Africa) the basal therocephalians Scylacosauridae and Lycosuchidae, gorgonopsians evolved from small and uncommon forms into large apex predators. Through the Middle to Upper Permian, in South Africa the dicynodonts were the most common animals, whereas the pareiasaurs Deltavjatia and Scutosaurus were the most abundant in the gorgonopsian-bearing Russian formations. During the Upper Permian, the South African Beaufort Group was a semi-arid cold steppe featuring large, seasonal (ephemeral) rivers and floodplains draining water sources much farther north into the Karoo Sea, with some occurrences of flash floods after sudden, heavy rainfall; the distribution of carbonates is consistent with present-day caliche deposits which form in climates with an average temperature of and of seasonal rainfall. The gorgonopsian-bearing Salarevskian Formation in western Russia was also probably deposited in a semi-arid environment with highly seasonal rainfall, and featured hygrophyte and halophyte plants in coastal areas, as well as more drought-resistant conifers at higher elevations. The Moradi Formation was an arid desert, primarily dominated by the captorhinid reptile Moradisaurus and the pareiasaur Bunostegos. It featured voltzian conifers, and has environmentally been compared to the interior Namib Desert or the Lake Eyre basin.
Extinction
Most gorgonopsians went extinct at the end of the Upper Permian during the Permian–Triassic extinction event, which was primarily caused by volcanism which formed the Siberian Traps. The resultant massive spike in greenhouse gases caused rapid aridification due to: temperature spike (as much as 8–10 °C at the equator, with average equatorial temperatures of 32–35 °C, or 90–95 °F, at the beginning of the Triassic), acid rain (with pH as low as 2 or 3 during eruption and 4 globally, and the subsequent dearth of forests for the first 10 million years of the Triassic), frequent wildfires (though they were already rather common throughout the Permian), and potential breakdown of the ozone layer (possibly briefly increasing UV radiation bombardment by 400% at the equator and 5000% at the poles). A possible specimen of Cyonosaurus'' suggests that some smaller taxa may have survived up to the latest Permian or Early Triassic. Among therapsids, small therocephalians and large herbivorous anomodonts managed to cross the Permian–Triassic boundary, and survived respectively until the Middle and Upper Triassic, but only small-bodied species of cynodonts survived into the Jurassic, whose descendants would include mammals. The niches gorgonopsians left open were eventually filled by the archosaurs (including crocodiles and dinosaurs) during the early stages of the Triassic.
| Biology and health sciences | Proto-mammals | Animals |
21183492 | https://en.wikipedia.org/wiki/Deep-sea%20community | Deep-sea community | A deep-sea community is any community of organisms associated by a shared habitat in the deep sea. Deep sea communities remain largely unexplored, due to the technological and logistical challenges and expense involved in visiting this remote biome. Because of the unique challenges (particularly the high barometric pressure, extremes of temperature, and absence of light), it was long believed that little life existed in this hostile environment. Since the 19th century however, research has demonstrated that significant biodiversity exists in the deep sea.
The three main sources of energy and nutrients for deep sea communities are marine snow, whale falls, and chemosynthesis at hydrothermal vents and cold seeps.
History
Prior to the 19th century scientists assumed life was sparse in the deep ocean. In the 1870s Sir Charles Wyville Thomson and colleagues aboard the Challenger expedition discovered many deep-sea creatures of widely varying types.
The first discovery of any deep-sea chemosynthetic community including higher animals was unexpectedly made at hydrothermal vents in the eastern Pacific Ocean during geological explorations (Corliss et al., 1979). Two scientists, J. Corliss and J. van Andel, first witnessed dense chemosynthetic clam beds from the submersible DSV Alvin on February 17, 1977, after their unanticipated discovery using a remote camera sled two days before.
The Challenger Deep is the deepest surveyed point of all of Earth's oceans; it is located at the southern end of the Mariana Trench near the Mariana Islands group. The depression is named after HMS Challenger, whose researchers made the first recordings of its depth on 23 March 1875 at station 225. The reported depth was 4,475 fathoms (8184 meters) based on two separate soundings. In 1960, Don Walsh and Jacques Piccard descended to the bottom of the Challenger Deep in the Trieste bathyscaphe. At this great depth a small flounder-like fish was seen moving away from the spotlight of the bathyscaphe.
The Japanese remote operated vehicle (ROV) Kaiko became the second vessel to reach the bottom of the Challenger Deep in March 1995. Nereus, a hybrid remotely operated vehicle (HROV) of the Woods Hole Oceanographic Institution, is the only vehicle capable of exploring ocean depths beyond 7000 meters. Nereus reached a depth of 10,902 meters at the Challenger Deep on May 31, 2009. On 1 June 2009, sonar mapping of the Challenger Deep by the Simrad EM120 multibeam sonar bathymetry system aboard the R/V Kilo Moana indicated a maximum depth of . The sonar system uses phase and amplitude bottom detection, with an accuracy of better than 0.2% of water depth (this is an error of about 22 meters at this depth).
Environment
Darkness
The ocean can be conceptualized as being divided into various zones, depending on depth, and presence or absence of sunlight. Nearly all life forms in the ocean depend on the photosynthetic activities of phytoplankton and other marine plants to convert carbon dioxide into organic carbon, which is the basic building block of organic matter. Photosynthesis in turn requires energy from sunlight to drive the chemical reactions that produce organic carbon.
The stratum of the water column up till which sunlight penetrates is referred to as the photic zone. The photic zone can be subdivided into two different vertical regions. The uppermost portion of the photic zone, where there is adequate light to support photosynthesis by phytoplankton and plants, is referred to as the euphotic zone (also referred to as the epipelagic zone, or surface zone). The lower portion of the photic zone, where the light intensity is insufficient for photosynthesis, is called the dysphotic zone (dysphotic means "poorly lit" in Greek). The dysphotic zone is also referred to as the mesopelagic zone, or the twilight zone. Its lowermost boundary is at a thermocline of , which, in the tropics generally lies between 200 and 1000 meters.
The euphotic zone is somewhat arbitrarily defined as extending from the surface to the depth where the light intensity is approximately 0.1–1% of surface sunlight irradiance, depending on season, latitude and degree of water turbidity. In the clearest ocean water, the euphotic zone may extend to a depth of about 150 meters, or rarely, up to 200 meters. Dissolved substances and solid particles absorb and scatter light, and in coastal regions the high concentration of these substances causes light to be attenuated rapidly with depth. In such areas the euphotic zone may be only a few tens of meters deep or less. The dysphotic zone, where light intensity is considerably less than 1% of surface irradiance, extends from the base of the euphotic zone to about 1000 meters. Extending from the bottom of the photic zone down to the seabed is the aphotic zone, a region of perpetual darkness.
Since the average depth of the ocean is about 3688 meters, the photic zone represents only a tiny fraction of the ocean's total volume. However, due to its capacity for photosynthesis, the photic zone has the greatest biodiversity and biomass of all oceanic zones. Nearly all primary production in the ocean occurs here. Any life forms present in the aphotic zone must either be capable of movement upwards through the water column into the photic zone for feeding, or must rely on material sinking from above, or must find another source of energy and nutrition, such as occurs in chemosynthetic archaea found near hydrothermal vents and cold seeps.
Hyperbaricity
These animals have evolved to survive the extreme pressure of the sub-photic zones. The pressure increases by about one atmosphere every ten meters. To cope with the pressure, many fish are rather small, usually not exceeding 25 cm in length. Also, scientists have discovered that the deeper these creatures live, the more gelatinous their flesh and more minimal their skeletal structure. These creatures have also eliminated all excess cavities that would collapse under the pressure, such as swim bladders.
Pressure is the greatest environmental factor acting on deep-sea organisms. In the deep sea, although most of the deep sea is under pressures between 200 and 600 atm, the range of pressure is from 20 to 1,000 atm. Pressure exhibits a great role in the distribution of deep sea organisms. Until recently, people lacked detailed information on the direct effects of pressure on most deep-sea organisms, because virtually all organisms trawled from the deep sea arrived at the surface dead or dying. With the advent of traps that incorporate a special pressure-maintaining chamber, undamaged larger metazoan animals have been retrieved from the deep sea in good condition. Some of these have been maintained for experimental purposes, and we are obtaining more knowledge of the biological effects of pressure.
Temperature
The two areas of greatest and most rapid temperature change in the oceans are the transition zone between the surface waters and the deep waters, the thermocline, and the transition between the deep-sea floor and the hot water flows at the hydrothermal vents. Thermoclines vary in thickness from a few hundred meters to nearly a thousand meters. Below the thermocline, the water mass of the deep ocean is cold and far more homogeneous. Thermoclines are strongest in the tropics, where the temperature of the epipelagic zone is usually above 20 °C. From the base of the epipelagic, the temperature drops over several hundred meters to 5 or 6 °C at 1,000 meters. It continues to decrease to the bottom, but the rate is much slower. Below 3,000 to 4,000 m, the water is isothermal. At any given depth, the temperature is practically unvarying over long periods of time. There are no seasonal temperature changes, nor are there any annual changes. No other habitat on earth has such a constant temperature.
Hydrothermal vents are the direct contrast with constant temperature. In these systems, the temperature of the water as it emerges from the "black smoker" chimneys may be as high as 400 °C (it is kept from boiling by the high hydrostatic pressure) while within a few meters it may be back down to 2–4 °C.
Salinity
Salinity is constant throughout the depths of the deep sea. There are two notable exceptions to this rule:
In the Mediterranean Sea, water loss through evaporation greatly exceeds input from precipitation and river runoff. Because of this, salinity in the Mediterranean is higher than in the Atlantic Ocean. Evaporation is especially high in its eastern half, causing the water level to decrease and salinity to increase in this area. The resulting pressure gradient pushes relatively cool, low-salinity water from the Atlantic Ocean across the basin. This water warms and becomes saltier as it travels eastward, then sinks in the region of the Levant and circulates westward, to spill back into the Atlantic over the Strait of Gibraltar. The net effect of this is that at the Strait of Gibraltar, there is an eastward surface current of cold water of lower salinity from the Atlantic, and a simultaneous westward current of warm saline water from the Mediterranean in the deeper zones. Once back in the Atlantic, this chemically distinct Mediterranean Intermediate Water can persist for thousands of kilometers away from its source.
Brine pools are large areas of brine on the seabed. These pools are bodies of water that have a salinity that is three to five times greater than that of the surrounding ocean. For deep sea brine pools the source of the salt is the dissolution of large salt deposits through salt tectonics. The brine often contains high concentrations of methane, providing energy to chemosynthetic extremophiles that live in this specialized biome. Brine pools are also known to exist on the Antarctic continental shelf where the source of brine is salt excluded during formation of sea ice. Deep sea and Antarctic brine pools can be toxic to marine animals. Brine pools are sometimes called seafloor lakes because the dense brine creates a halocline which does not easily mix with overlying seawater. The high salinity raises the density of the brine, which creates a distinct surface and shoreline for the pool.
The deep sea, or deep layer, is the lowest layer in the ocean, existing below the thermocline, at a depth of 1000 fathoms (1800 m) or more. The deepest part of the deep sea is Mariana Trench located in the western North Pacific. It is also the deepest point of the Earth's crust. It has a maximum depth of about 10.9 km which is deeper than the height of Mount Everest. In 1960, Don Walsh and Jacques Piccard reached the bottom of Mariana Trench in the Trieste bathyscaphe. The pressure is about 11,318 metric tons-force per square meter (110.99 MPa or 16100 psi).
Zones
Mesopelagic
The mesopelagic zone is the upper section of the midwater zone, and extends from below sea level. This is colloquially known as the "twilight zone" as light can still penetrate this layer, but it is too low to support photosynthesis. The limited amount of light, however, can still allow organisms to see, and creatures with a sensitive vision can detect prey, communicate, and orientate themselves using their sight. Organisms in this layer have large eyes to maximize the amount of light in the environment.
Most mesopelagic fish make daily vertical migrations, moving at night into the epipelagic zone, often following similar migrations of zooplankton, and returning to the depths for safety during the day. These vertical migrations often occur over a large vertical distances, and are undertaken with the assistance of a swimbladder. The swimbladder is inflated when the fish wants to move up, and, given the high pressures in the mesopelagic zone, this requires significant energy. As the fish ascends, the pressure in the swimbladder must adjust to prevent it from bursting. When the fish wants to return to the depths, the swimbladder is deflated. Some mesopelagic fishes make daily migrations through the thermocline, where the temperature changes between , thus displaying considerable tolerances for temperature change.
Mesopelagic fish usually lack defensive spines, and use colour and bioluminescence to camouflage them from other fish. Ambush predators are dark, black or red. Since the longer, red, wavelengths of light do not reach the deep sea, red effectively functions the same as black. Migratory forms use countershaded silvery colours. On their bellies, they often display photophores producing low grade light. For a predator from below, looking upwards, this bioluminescence camouflages the silhouette of the fish. However, some of these predators have yellow lenses that filter the (red deficient) ambient light, leaving the bioluminescence visible.
Bathyal
The bathyal zone is the lower section of the midwater zone, and encompasses the depths of . Light does not reach this zone, giving it its nickname "the midnight zone"; due to the lack of light, it is less densely populated than the epipelagic zone, despite being much larger. Fish find it hard to live in this zone, as there is crushing pressure, cold temperatures of , a low level of dissolved oxygen, and a lack of sufficient nutrients. What little energy is available in the bathypelagic zone filters from above in the form of detritus, faecal material, and the occasional invertebrate or mesopelagic fish. About 20% of the food that has its origins in the epipelagic zone falls down to the mesopelagic zone, but only about 5% filters down to the bathypelagic zone. The fish that do live there may have reduced or completely lost their gills, kidneys, hearts, and swimbladders, have slimy instead of scaly skin, and have a weak skeletal and muscular build. This lack of ossification is an adaptation to save energy when food is scarce.
Most of the animals that live in the bathyal zone are invertebrates such as sea sponges, cephalopods, and echinoderms. With the exception of very deep areas of the ocean, the bathyal zone usually reaches the benthic zone on the seafloor.
Abyssal and hadal
The abyssal zone remains in perpetual darkness at a depth of . The only organisms that inhabit this zone are chemotrophs and predators that can withstand immense pressures, sometimes as high as . The hadal zone (named after Hades, the Greek god of the underworld) is a zone designated for the deepest trenches in the world, reaching depths of below . The deepest point in the hadal zone is the Marianas Trench, which descends to and has a pressure of .
Energy sources
Marine snow
The upper photic zone of the ocean is filled with particulate organic matter (POM), especially in the coastal areas and the upwelling areas. However, most POM is small and light. It may take hundreds, or even thousands of years for these particles to settle through the water column into the deep ocean. This time delay is long enough for the particles to be remineralized and taken up by organisms in the food web.
In the deep Sargasso Sea, scientists from the Woods Hole Oceanographic Institution (WHOI) found what became known as marine snow in which the POM are repackaged into much larger particles which sink at much greater speed, falling like snow.
Because of the sparsity of food, the organisms living on and in the bottom are generally opportunistic. They have special adaptations for this extreme environment: rapid growth, effect larval dispersal mechanism and the ability to use a 'transient' food resource. One typical example is wood-boring bivalves, which bore into wood and other plant remains and are fed on the organic matter from the remains.
Occasional surface blooms
Sometimes sudden access to nutrients near the surface leads to blooms of plankton, algae or animals such as salps, which becomes so numerous that they will sink all the way to the bottom without being consumed by other organisms. These short bursts of nutrients reaching the seafloor can exceed years of marine snow, and are rapidly consumed by animals and microbes. The waste products becomes part of the deep-sea sediments, and recycled by animals and microbes that feed on mud for years to come.
Whale falls
For the deep-sea ecosystem, the death of a whale is the most important event. A dead whale can bring hundreds of tons of organic matter to the bottom. Whale fall community progresses through three stages:
Mobile scavenger stage: Big and mobile deep-sea animals arrive at the site almost immediately after whales fall on the bottom. Amphipods, crabs, sleeper sharks and hagfish are all scavengers.
Opportunistic stage: Organisms arrive which colonize the bones and surrounding sediments that have been contaminated with organic matter from the carcass and any other tissue left by the scavengers. One genus is Osedax, a tube worm. The larva is born without sex. The surrounding environment determines the sex of the larva. When a larva settles on a whale bone, it turns into a female; when a larva settles on or in a female, it turns into a dwarf male. One female Osedax can carry more than 200 of these male individuals in its oviduct.
Sulfophilic stage: Further decomposition of bones and seawater sulfate reduction happen at this stage. Bacteria create a sulphide-rich environment analogous to hydrothermal vents. Polynoids, bivalves, gastropods and other sulphur-loving creatures move in.
Chemosynthesis
Hydrothermal vents
Hydrothermal vents were discovered in 1977 by scientists from Scripps Institution of Oceanography. So far, the discovered hydrothermal vents are all located at the boundaries of plates: East Pacific, California, Mid-Atlantic ridge, China and Japan.
New ocean basin material is being made in regions such as the Mid-Atlantic ridge as tectonic plates pull away from each other. The rate of spreading of plates is 1–5 cm/yr. Cold sea water circulates down through cracks between two plates and heats up as it passes through hot rock. Minerals and sulfides are dissolved into the water during the interaction with rock. Eventually, the hot solutions emanate from an active sub-seafloor rift, creating a hydrothermal vent.
Chemosynthesis of bacteria provide the energy and organic matter for the whole food web in vent ecosystems. These vents spew forth very large amounts of chemicals, which these bacteria can transform into energy. These bacteria can also grow free of a host and create mats of bacteria on the sea floor around hydrothermal vents, where they serve as food for other creatures. Bacteria are a key energy source in the food chain. This source of energy creates large populations in areas around hydrothermal vents, which provides scientists with an easy stop for research. Organisms can also use chemosynthesis to attract prey or to attract a mate.
Giant tube worms can grow to tall because of the richness of nutrients. Over 300 new species have been discovered at hydrothermal vents.
Hydrothermal vents are entire ecosystems independent from sunlight, and may be the first evidence that the earth can support life without the sun.
Cold seeps
A cold seep (sometimes called a cold vent) is an area of the ocean floor where hydrogen sulfide, methane and other hydrocarbon-rich fluid seepage occurs, often in the form of a brine pool.
Ecology
Deep sea food webs are complex, and aspects of the system are poorly understood. Typically, predator-prey interactions within the deep are compiled by direct observation (likely from remotely operated underwater vehicles), analysis of stomach contents, and biochemical analysis. Stomach content analysis is the most common method used, but it is not reliable for some species.
In deep sea pelagic ecosystems off of California, the trophic web is dominated by deep sea fishes, cephalopods, gelatinous zooplankton, and crustaceans. Between 1991 and 2016, 242 unique feeding relationships between 166 species of predators and prey demonstrated that gelatinous zooplankton have an ecological impact similar to that of large fishes and squid. Narcomedusae, siphonophores (of the family Physonectae), ctenophores, and cephalopods consumed the greatest diversity of prey, in decreasing order. Cannibalism has been documented in squid of the genus Gonatus.
Deep sea mining has severe consequences for ocean ecosystems. The destruction of habitats, disturbance of sediment layers, and noise pollution threaten marine species. Essential biodiversity can be lost, with unpredictable effects on the food chain. Additionally, toxic metals and chemicals can be released, leading to long-term pollution of seawater. This raises questions about the sustainability and environmental costs of such activities.
Deep sea research
Humans have explored less than 4% of the ocean floor, and dozens of new species of deep sea creatures are discovered with every dive. The submarine DSV Alvin—owned by the US Navy and operated by the Woods Hole Oceanographic Institution (WHOI) in Woods Hole, Massachusetts—exemplifies the type of craft used to explore deep water. This 16 ton submarine can withstand extreme pressure and is easily manoeuvrable despite its weight and size.
The extreme difference in pressure between the sea floor and the surface makes creatures' survival on the surface near impossible; this makes in-depth research difficult because most useful information can only be found while the creatures are alive. Recent developments have allowed scientists to look at these creatures more closely, and for a longer time. Marine biologist Jeffery Drazen has explored a solution: a pressurized fish trap. This captures a deep-water creature, and adjusts its internal pressure slowly to surface level as the creature is brought to the surface, in the hope that the creature can adjust.
Another scientific team, from the Université Pierre-et-Marie-Curie, has developed a capture device known as the PERISCOP, which maintains water pressure as it surfaces, thus keeping the samples in a pressurized environment during the ascent. This permits close study on the surface without any pressure disturbances affecting the sample.
| Physical sciences | Water: General | Earth science |
21188370 | https://en.wikipedia.org/wiki/Fuel | Fuel | A fuel is any material that can be made to react with other substances so that it releases energy as thermal energy or to be used for work. The concept was originally applied solely to those materials capable of releasing chemical energy but has since also been applied to other sources of heat energy, such as nuclear energy (via nuclear fission and nuclear fusion).
The heat energy released by reactions of fuels can be converted into mechanical energy via a heat engine. Other times, the heat itself is valued for warmth, cooking, or industrial processes, as well as the illumination that accompanies combustion. Fuels are also used in the cells of organisms in a process known as cellular respiration, where organic molecules are oxidized to release usable energy. Hydrocarbons and related organic molecules are by far the most common source of fuel used by humans, but other substances, including radioactive metals, are also utilized.
Fuels are contrasted with other substances or devices storing potential energy, such as those that directly release electrical energy (such as batteries and capacitors) or mechanical energy (such as flywheels, springs, compressed air, or water in a reservoir).
History
The first known use of fuel was the combustion of firewood by Homo erectus nearly two million years ago. Throughout most of human history only fuels derived from plants or animal fat were used by humans. Charcoal, a wood derivative, has been used since at least 6,000 BCE for melting metals. It was only supplanted by coke, derived from coal, as European forests started to become depleted around the 18th century. Charcoal briquettes are now commonly used as a fuel for barbecue cooking.
Crude oil was distilled by Persian chemists, with clear descriptions given in Arabic handbooks such as those of Muhammad ibn Zakarīya Rāzi. He described the process of distilling crude oil/petroleum into kerosene, as well as other hydrocarbon compounds, in his Kitab al-Asrar (Book of Secrets). Kerosene was also produced during the same period from oil shale and bitumen by heating the rock to extract the oil, which was then distilled. Rāzi also gave the first description of a kerosene lamp using crude mineral oil, referring to it as the "naffatah".
The streets of Baghdad were paved with tar, derived from petroleum that became accessible from natural fields in the region. In the 9th century, oil fields were exploited in the area around modern Baku, Azerbaijan. These fields were described by the Arab geographer Abu al-Hasan 'Alī al-Mas'ūdī in the 10th century, and by Marco Polo in the 13th century, who described the output of those wells as hundreds of shiploads.
With the development of the steam engine in the United Kingdom in 1769, coal came into more common use, the combustion of which releases chemical energy that can be used to turn water into steam. Coal was later used to drive ships and locomotives. By the 19th century, gas extracted from coal was being used for street lighting in London. In the 20th and 21st centuries, the primary use of coal is to generate electricity, providing 40% of the world's electrical power supply in 2005.
Fossil fuels were rapidly adopted during the Industrial Revolution, because they were more concentrated and flexible than traditional energy sources, such as water power. They have become a pivotal part of our contemporary society, with most countries in the world burning fossil fuels in order to produce power, but are falling out of favor due to the global warming and related effects that are caused by burning them.
Currently the trend has been towards renewable fuels, such as biofuels like alcohols.
Chemical
Chemical fuels are substances that release energy by reacting with substances around them, most notably by the process of combustion.
Chemical fuels are divided in two ways. First, by their physical properties, as a solid, liquid or gas. Secondly, on the basis of their occurrence: primary (natural fuel) and secondary (artificial fuel). Thus, a general classification of chemical fuels is:
Solid fuel
Solid fuel refers to various types of solid material that are used as fuel to produce energy and provide heating, usually released through combustion. Solid fuels include wood, charcoal, peat, coal, hexamine fuel tablets, and pellets made from wood (see wood pellets), corn, wheat, rye and other grains. Solid-fuel rocket technology also uses solid fuel (see solid propellants). Solid fuels have been used by humanity for many years to create fire. Coal was the fuel source which enabled the Industrial Revolution, from firing furnaces, to running steam engines. Wood was also extensively used to run steam locomotives. Both peat and coal are still used in electricity generation today.
The use of some solid fuels (e.g. coal) is restricted or prohibited in some urban areas, due to unsafe levels of toxic emissions. The use of other solid fuels as wood is decreasing as heating technology and the availability of good quality fuel improves. In some areas, smokeless coal is often the only solid fuel used. In Ireland, peat briquettes are used as smokeless fuel. They are also used to start a coal fire.
Liquid fuels
Liquid fuels are combustible or energy-generating molecules that can be harnessed to create mechanical energy, usually producing kinetic energy. They must also take the shape of their container; the fumes of liquid fuels are flammable, not the fluids.
Most liquid fuels in widespread use are derived from the fossilized remains of dead plants and animals by exposure to heat and pressure inside the Earth's crust. However, there are several types, such as hydrogen fuel (for automotive uses), ethanol, jet fuel and bio-diesel, which are all categorized as liquid fuels. Emulsified fuels of oil in water, such as orimulsion, have been developed as a way to make heavy oil fractions usable as liquid fuels. Many liquid fuels play a primary role in transportation and the economy.
Some common properties of liquid fuels are that they are easy to transport and can be handled easily. They are also relatively easy to use for all engineering applications and in home use. Fuels like kerosene are rationed in some countries, for example in government-subsidized shops in India for home use.
Conventional diesel is similar to gasoline in that it is a mixture of aliphatic hydrocarbons extracted from petroleum. Kerosene is used in kerosene lamps and as a fuel for cooking, heating, and small engines. Natural gas, composed chiefly of methane, can only exist as a liquid at very low temperatures (regardless of pressure), which limits its direct use as a liquid fuel in most applications. LP gas is a mixture of propane and butane, both of which are easily compressible gases under standard atmospheric conditions. It offers many of the advantages of compressed natural gas (CNG) but is denser than air, does not burn as cleanly, and is much more easily compressed. Commonly used for cooking and space heating, LP gas and compressed propane are seeing increased use in motorized vehicles. Propane is the third most commonly used motor fuel globally.
Fuel gas
Fuel gas is any one of a number of fuels that are gaseous under ordinary conditions. Many fuel gases are composed of hydrocarbons (such as methane or propane), hydrogen, carbon monoxide, or mixtures thereof. Such gases are sources of potential heat energy or light energy that can be readily transmitted and distributed through pipes from the point of origin directly to the place of consumption. Fuel gas is contrasted with liquid fuels and from solid fuels, though some fuel gases are liquefied for storage or transport. While their gaseous nature can be advantageous, avoiding the difficulty of transporting solid fuel and the dangers of spillage inherent in liquid fuels, it can also be dangerous. It is possible for a fuel gas to be undetected and collect in certain areas, leading to the risk of a gas explosion. For this reason, odorizers are added to most fuel gases so that they may be detected by a distinct smell. The most common type of fuel gas in current use is natural gas.
Biofuels
Biofuel can be broadly defined as solid, liquid, or gas fuel consisting of, or derived from biomass. Biomass can also be used directly for heating or power—known as biomass fuel. Biofuel can be produced from any carbon source that can be replenished rapidly e.g. plants. Many different plants and plant-derived materials are used for biofuel manufacture.
Perhaps the earliest fuel employed by humans is wood. Evidence shows controlled fire was used up to 1.5 million years ago at Swartkrans, South Africa. It is unknown which hominid species first used fire, as both Australopithecus and an early species of Homo were present at the sites. As a fuel, wood has remained in use up until the present day, although it has been superseded for many purposes by other sources. Wood has an energy density of 10–20 MJ/kg.
Recently biofuels have been developed for use in automotive transport (for example bioethanol and biodiesel), but there is widespread public debate about how carbon neutral these fuels are.
Fossil fuels
Fossil fuels are hydrocarbons, primarily coal and petroleum (liquid petroleum or natural gas), formed from the fossilized remains of ancient plants and animals by exposure to high heat and pressure in the absence of oxygen in the Earth's crust over hundreds of millions of years. Commonly, the term fossil fuel also includes hydrocarbon-containing natural resources that are not derived entirely from biological sources, such as tar sands. These latter sources are properly known as mineral fuels.
Fossil fuels contain high percentages of carbon and include coal, petroleum, and natural gas.
They range from volatile materials with low carbon:hydrogen ratios like methane, to liquid petroleum to nonvolatile materials composed of almost pure carbon, like anthracite coal. Methane can be found in hydrocarbon fields, alone, associated with oil, or in the form of methane clathrates. Fossil fuels formed from the fossilized remains of dead plants by exposure to heat and pressure in the Earth's crust over millions of years. This biogenic theory was first introduced by German scholar Georg Agricola in 1556 and later by Mikhail Lomonosov in the 18th century.
It was estimated by the Energy Information Administration that in 2007 primary sources of energy consisted of petroleum 36.0%, coal 27.4%, natural gas 23.0%, amounting to an 86.4% share for fossil fuels in primary energy consumption in the world. Non-fossil sources in 2006 included hydroelectric 6.3%, nuclear 8.5%, and others (geothermal, solar, tidal, wind, wood, waste) amounting to 0.9%. World energy consumption was growing about 2.3% per year.
Fossil fuels are non-renewable resources because they take millions of years to form, and reserves are being depleted much faster than new ones are being made. So we must conserve these fuels and use them judiciously. The production and use of fossil fuels raise environmental concerns. A global movement toward the generation of renewable energy is therefore under way to help meet increased energy needs. The burning of fossil fuels produces around 21.3 billion tonnes (21.3 gigatonnes) of carbon dioxide (CO2) per year, but it is estimated that natural processes can only absorb about half of that amount, so there is a net increase of 10.65 billion tonnes of atmospheric carbon dioxide per year (one tonne of atmospheric carbon is equivalent to (this is the ratio of the molecular/atomic weights) or 3.7 tonnes of CO2. Carbon dioxide is one of the greenhouse gases that enhances radiative forcing and contributes to global warming, causing the average surface temperature of the Earth to rise in response, which the vast majority of climate scientists agree will cause major adverse effects.
Fuels are a source of energy.
The International Energy Agency (IEA) predicts that fossil fuel prices will decline, with oil stabilizing around $75 to $80 per barrel as electric vehicle adoption surges and renewable energy expands. Additionally, the IEA anticipates a notable increase in liquefied natural gas capacity, enhancing Europe’s energy diversification.
Energy
The amount of energy from different types of fuel depends on the stoichiometric ratio, the chemically correct air and fuel ratio to ensure complete combustion of fuel, and its specific energy, the energy per unit mass.
| Technology | Energy and fuel | null |
6595512 | https://en.wikipedia.org/wiki/Lucy%20%28Australopithecus%29 | Lucy (Australopithecus) | AL 288-1, commonly known as Lucy or Dinkinesh (), is a collection of several hundred pieces of fossilized bone comprising 40 percent of the skeleton of a female of the hominin species Australopithecus afarensis. It was discovered in 1974 in Ethiopia, at Hadar, a site in the Awash Valley of the Afar Triangle, by Donald Johanson, a paleoanthropologist of the Cleveland Museum of Natural History.
Lucy is an early australopithecine and is dated to about 3.2 million years ago. The skeleton presents a small skull akin to that of non-hominin apes, plus evidence of a walking-gait that was bipedal and upright, akin to that of humans (and other hominins); this combination supports the view of human evolution that bipedalism preceded increase in brain size. A 2016 study proposes that Australopithecus afarensis was at least partly, tree-dwelling, though the extent of this is debated.
Lucy was named by Pamela Alderman after the 1967 song "Lucy in the Sky with Diamonds" by the Beatles, which was played loudly and repeatedly in the expedition camp all evening after the excavation team's first day of work on the recovery site. After public announcement of the discovery, Lucy captured much international interest, becoming a household name at the time.
Lucy became famous worldwide, and the story of her discovery and reconstruction was published in a book by Johanson and Edey. Beginning in 2007, the fossil assembly and associated artefacts were exhibited publicly in an extended six-year tour of the United States; the exhibition was called Lucy's Legacy: The Hidden Treasures of Ethiopia. There was discussion of the risks of damage to the unique fossils, and other museums preferred to display casts of the fossil assembly. The original fossils were returned to Ethiopia in 2013, and subsequent exhibitions have used casts.
Recent research has revealed that she is no longer considered the earliest known member of the human family. Contrary to earlier beliefs that her species first walked upright in open savanna grasslands, new evidence suggests they walked in grassy woodlands with deciduous trees. Her species adapted to various habitats over millennia, enduring changes in climate. Importantly, she was not alone in her environment. "We have multiple [hominin] species in the same time period," said Yohannes Haile-Selassie, director of the Institute of Human Origins at Arizona State University.
Discovery
Background controversy
Ever since the development of evolutionary theory in the early 19th century, biologists recognized that humans must be distantly related to all other species. Without transitional fossils, scientists presumed that humans' closest relatives were the great apes. They also assumed the first traits to evolve after speciation related to intelligence: big brains, tool use, and complex language.
In the 1920s, Raymond Dart discovered the Taung child. That skeleton seemed bipedal (unlike chimps), but lacked skull space for a powerful brain. Without further data to contextualize Dart's find, anthropologists could not prove whether bipedality, intelligence, or some other trait had first distinguished proto-humans from their great ape relatives.
Organizing the expedition
French geologist and paleoanthropologist Maurice Taieb discovered the Hadar Formation for paleoanthropology in 1970 in the Afar Triangle of Ethiopia, then in Hararghe province; he recognized its potential as a likely repository of the fossils and artifacts of human origins. Taieb formed the International Afar Research Expedition (IARE) and invited three prominent international scientists to conduct research expeditions into the region. Under his directorship, these were: Donald Johanson (co-director), an American paleoanthropologist and curator at the Cleveland Museum of Natural History, who later founded the Institute of Human Origins, now part of Arizona State University; Yves Coppens (1934–2022, co-director), a French paleoanthropologist appointed in 1983 a professor at the Collège de France, which is considered to be France's most prestigious research establishment, and Mary Leakey, the noted British paleoanthropologist. An expedition was soon mounted with seven French and four American participants; in the autumn of 1973 the team began surveying sites around Hadar for signs related to the origin of humans.
First find
In November 1973, near the end of the first field season, Johanson noticed a fossil of the upper end of a shinbone, which had been sliced slightly at the front. The lower end of a femur was found near it, and when he fitted them together, the angle of the knee joint clearly showed that this fossil, reference AL 129-1, was an upright walking hominin. This fossil was later dated at more than three million years old—much older than other hominin fossils known at the time. The site lay about from the site where "Lucy" subsequently was found, in a rock stratum deeper than that in which the Lucy fragments were found.
Subsequent findings
The team returned for the second field season the following year and found hominin jaws. Then, on the morning of November 24, 1974, near the Awash River, Johanson abandoned a plan to update his field notes and joined graduate student Tom Gray to search Locality 162 for bone fossils.
By Johanson's later (published) accounts, both he and Tom Gray spent two hours on the increasingly hot and arid plain, surveying the dusty terrain. On a hunch, Johanson decided to look at the bottom of a small gully that had been checked at least twice before by other workers. At first view nothing was immediately visible, but as they turned to leave a fossil caught Johanson's eye; an arm bone fragment was lying on the slope. Near it lay a fragment from the back of a small skull. They noticed part of a femur (thigh bone) a few feet (about one meter) away. As they explored further, they found more and more bones on the slope, including vertebrae, part of a pelvis, ribs, and pieces of jaw. They marked the spot and returned to camp, excited at finding so many pieces apparently from one individual hominin.
In the afternoon, all members of the expedition returned to the gully to section off the site and prepare it for careful excavation and collection, which eventually took three weeks. That first evening they celebrated at the camp; at some stage during the evening they named fossil AL 288-1 "Lucy", after the Beatles' song "Lucy in the Sky with Diamonds" (1967), which was being played loudly and repeatedly on a tape recorder in the camp.
Over the next three weeks the team found several hundred pieces or fragments of bone with no duplication, confirming their original speculation that the pieces were from a single individual; ultimately, it was determined that an amazing 40 percent of a hominin skeleton was recovered at the site. Johanson assessed it as female based on the one complete pelvic bone and sacrum, which indicated the width of the pelvic opening.
Assembling the pieces
Lucy was tall, weighed , and (after reconstruction) looked somewhat like a chimpanzee. The creature had a small brain like a chimpanzee, but the pelvis and leg bones were almost identical in function to those of modern humans, showing with certainty that Lucy's species were hominins that had stood upright and had walked erect.
Reconstruction in Cleveland
With the permission of the government of Ethiopia, Johanson brought all the skeletal fragments to the Cleveland Museum of Natural History in Ohio, where they were stabilized and reconstructed by anthropologist Owen Lovejoy. Lucy the pre-human hominid and fossil hominin, captured much public notice; she became almost a household name at the time. Some nine years later, and fully assembled, she was returned to Ethiopia.
Later discoveries
Additional finds of A. afarensis were made during the 1970s and forward, gaining for anthropologists a better understanding of the ranges of morphic variability and sexual dimorphism within the species. A more complete skeleton of a related hominid, Ardipithecus, was found in the same Awash Valley in 1992. "Ardi", like "Lucy", was a hominid-becoming-hominin species, but, dated at , it had evolved much earlier than the afarensis species. Excavation, preservation, and analysis of the specimen Ardi was very difficult and time-consuming; work was begun in 1992, with the results not fully published until October 2009.
Fossil age estimates
Initial attempts were made in 1974 by Maurice Taieb and James Aronson in Aronson's laboratory at Case Western Reserve University to estimate the age of the fossils using the potassium-argon radiometric dating method. These efforts were hindered by several factors: the rocks in the recovery area were chemically altered or reworked by volcanic activity; datable crystals were very scarce in the sample material; and there was a complete absence of pumice clasts at Hadar. (The Lucy skeleton occurs in the part of the Hadar sequence that accumulated with the fastest rate of deposition, which partly accounts for her excellent preservation.)
Fieldwork at Hadar was suspended in the winter of 1976–77. When it was resumed thirteen years later in 1990, the more precise argon-argon technology had been updated by Derek York at the University of Toronto. By 1992 Aronson and Robert Walter had found two suitable samples of volcanic ash—the older layer of ash was about 18 m below the fossil and the younger layer was only one meter below, closely marking the age of deposition of the specimen. These samples were argon-argon dated by Walter in the geochronology laboratory of the Institute of Human Origins at 3.22 and 3.18 million years.
Skeletal characteristics
Ambulation
One of the most striking characteristics of the Lucy skeleton is a valgus knee, which indicates that she normally moved by walking upright. Her femur presents a mix of ancestral and derived traits. The femoral head is small and the femoral neck is short; both are primitive traits. The greater trochanter, however, is clearly a derived trait, being short and human-like—even though, unlike in humans, it is situated higher than the femoral head. The length ratio of her humerus (arm) to femur (thigh) is 84.6%, which compares to 71.8% for modern humans, and 97.8% for common chimpanzees, indicating that either the arms of A. afarensis were beginning to shorten, the legs were beginning to lengthen, or both were occurring simultaneously. Lucy also had a lordose curve, or lumbar curve, another indicator of habitual bipedalism. She apparently had physiological flat feet, not to be confused with pes planus or any pathology, even though other afarensis individuals appear to have had arched feet.
Pelvic girdle
Johanson recovered Lucy's left innominate bone and sacrum. Though the sacrum was remarkably well preserved, the innominate was distorted, leading to two different reconstructions. The first reconstruction had little iliac flare and virtually no anterior wrap, creating an ilium that greatly resembled that of an ape. However, this reconstruction proved to be faulty, as the superior pubic rami would not have been able to connect were the right ilium identical to the left.
A later reconstruction by Tim White showed a broad iliac flare and a definite anterior wrap, indicating that Lucy had an unusually broad inner acetabular distance and unusually long superior pubic rami. Her pubic arch was over 90 degrees and derived; that is, similar to modern human females. Her acetabulum, however, was small and primitive.
Sacrum and spine
While examining Lucy's fossilized remains, it was believed that Lucy's sacrum had five fused elements. The sacrum was found to be in good condition with little damage done. Although the sacrum had five fused elements, the transverse processes of the most caudal element were not seen to connect to the segments craniad to it. This would result in researchers concluding that the sacrum suffered fossil damage which led to the fifth segment not connecting.
Although this was the case, in the mid-2010s studies came out with new theories as to why Lucy's fifth sacral segment is in that shape. Some researchers conclude that Lucy has only four sacral segments. Published in the American Journal of Physical Anthropology, researchers suggest that fossil damage did not shorten the transverse process and that Lucy's sacrum was in this state from the beginning. This specific study points to Lucy's sacrum having four sacral segments which researchers say conforms with the "long-back" model of hominoid vertebral evolution. There are some disagreements in the community about the fifth sacral segment and if fossil damage was enough to change the fifth segment or if it was originally in that state. Discussed in the Journal, researchers conclude that Lucy having only four sacral segments is consistent with other findings related to early Miocene hominoids.
Lucy's back is associated with approximately 9 vertebrae. Although Lucy was found with a relatively intact and well preserved sacrum, she was missing pieces in her spinal column. Lucy's discoverers and later workers had given the vertebrae provisional level assignments to locations within the vertebral column. Some vertebrae were in worse conditions than others. Lucy had a worn out upper thoracic neutral arch. Researchers have yet to find a cause as to why this particular vertebrae was in worse condition than the other pieces. While accessing and restructuring Lucy's spinal column, it was noted to have been missing pieces that leave it incomplete. Not including an oddly worn out upper thoracic neutral arch, and the lumbar vertebrae, the other remaining thoracic vertebrae were compiled to form an incomplete formation. The formation was arranged from the sixth thoracic vertebrae (T6) to its caudal end (T12), with the seventh thoracic vertebrae (T7) missing. the continuity differs in the thoracic series between researchers and is being reevaluated. Although there are new studies and reassessments being done, it does not refute previous work or conclusions about Lucy's spine.
Cranial specimens
The cranial evidence recovered from Lucy is far less derived than her postcranium. Her neurocranium is small and primitive, while she possesses more spatulate canines than other apes. The cranial capacity was about 375 to 500 cubic centimeters.
Rib cage and plant-based diet
Australopithecus afarensis seems to have had the same conical rib-cage found in today's non-human great apes (like the chimpanzee and gorilla), which allows room for a large stomach and the longer intestine needed for digesting voluminous plant matter. Fully 60% of the blood supply of non-human apes is used in the digestion process, greatly impeding the development of brain function (which is limited thereby to using about 10% of the circulation). The heavier musculature of the jaws—those muscles operating the intensive masticatory process for chewing plant material—similarly would also limit development of the braincase. During evolution of the human lineage these muscles seem to have weakened with the loss of the myosin gene MYH16, a two base-pair deletion that occurred possibly about 2.4 million years ago.
Other findings
A study of the mandible across a number of specimens of A. afarensis indicated that Lucy's jaw was rather unlike other hominins, having a more gorilla-like appearance. Rak et al. concluded that this morphology arose "independently in gorillas and hominins", and that A. afarensis is "too derived to occupy a position as a common ancestor of both the Homo and robust australopith clades".
Work at the American Museum of Natural History uncovered a possible Theropithecus vertebral fragment that was found mixed in with Lucy's vertebrae, but confirmed the remainder belonged to her.
Death
Lucy's cause of death has not been determined. The specimen does not show the signs of post-mortem bone damage characteristic of animals killed by predators and then scavenged. The only visible damage is a single carnivore tooth mark on the top of her left pubic bone, believed to have occurred at or around the time of death, but which is not necessarily related to her death. Her third molars were erupted and slightly worn and, therefore, it was concluded that she was fully matured with completed skeletal development. There are indications of degenerative disease to her vertebrae that do not necessarily indicate old age. It is believed that she was a mature but young adult when she died.
In 2016 researchers at the University of Texas at Austin suggested that Lucy died after falling from a tall tree. Donald Johanson and Tim White disagreed with the suggestions.
Exhibitions
The Lucy skeleton is preserved at the National Museum of Ethiopia in Addis Ababa. A plaster replica is publicly displayed there instead of the original skeleton. A cast of the original skeleton in its reconstructed form is displayed at the Cleveland Museum of Natural History. At the American Museum of Natural History in New York City a diorama presents Australopithecus afarensis and other human predecessors, showing each species and its habitat and explaining the behaviors and capabilities assigned to each. A cast of the skeleton as well as a corpus reconstruction of Lucy is displayed at The Field Museum in Chicago.
US tour
A six-year exhibition tour of the United States was undertaken during 2007–13; it was titled Lucy's Legacy: The Hidden Treasures of Ethiopia and it featured the actual Lucy fossil reconstruction and over 100 artifacts from prehistoric times to the present. The tour was organized by the Houston Museum of Natural Science and was approved by the Ethiopian government and the U.S. State Department. A portion of the proceeds from the tour was designated to modernizing Ethiopia's museums.
There was controversy in advance of the tour over concerns about the fragility of the specimens, with various experts including paleoanthropologist Owen Lovejoy and anthropologist and conservationist Richard Leakey publicly stating their opposition, while discoverer Don Johanson, despite concerns for the possibility of damage, felt the tour would raise awareness of human origins studies. The Smithsonian Institution, Cleveland Museum of Natural History and other museums declined to host the exhibits.
The Houston Museum made arrangements for exhibiting at ten other museums, including the Pacific Science Center in Seattle. In September 2008, between the exhibits in Houston and Seattle, the skeletal assembly was taken to the University of Texas at Austin for 10 days to perform high-resolution CT scans of the fossils.
Lucy was exhibited at the Discovery Times Square Exposition in New York City from June until October 2009. In New York, the exhibition included Ida (Plate B), the other half of the recently announced Darwinius masilae fossil. She was also exhibited in Mexico at the Mexico Museum of Anthropology until its return to Ethiopia in May 2013.
Ethiopia celebrated the return of Lucy in May 2013.
| Biology and health sciences | Australopithecines | Biology |
18914017 | https://en.wikipedia.org/wiki/Alzheimer%27s%20disease | Alzheimer's disease | Alzheimer's disease (AD) is a neurodegenerative disease that usually starts slowly and progressively worsens. It is the cause of 60–70% of cases of dementia. The most common early symptom is difficulty in remembering recent events. As the disease advances, symptoms can include problems with language, disorientation (including easily getting lost), mood swings, loss of motivation, self-neglect, and behavioral issues. As a person's condition declines, they often withdraw from family and society. Gradually, bodily functions are lost, ultimately leading to death. Although the speed of progression can vary, the average life expectancy following diagnosis is three to twelve years.
The causes of Alzheimer's disease remain poorly understood. There are many environmental and genetic risk factors associated with its development. The strongest genetic risk factor is from an allele of apolipoprotein E. Other risk factors include a history of head injury, clinical depression, and high blood pressure. The progression of the disease is largely characterized by the accumulation of malformed protein deposits in the cerebral cortex, called amyloid plaques and neurofibrillary tangles. These misfolded protein aggregates interfere with normal cell function, and over time lead to irreversible degeneration of neurons and loss of synaptic connections in the brain. A probable diagnosis is based on the history of the illness and cognitive testing, with medical imaging and blood tests to rule out other possible causes. Initial symptoms are often mistaken for normal brain aging. Examination of brain tissue is needed for a definite diagnosis, but this can only take place after death.
No treatments can stop or reverse its progression, though some may temporarily improve symptoms. A healthy diet, physical activity, and social engagement are generally beneficial in aging, and may help in reducing the risk of cognitive decline and Alzheimer's. Affected people become increasingly reliant on others for assistance, often placing a burden on caregivers. The pressures can include social, psychological, physical, and economic elements. Exercise programs may be beneficial with respect to activities of daily living and can potentially improve outcomes. Behavioral problems or psychosis due to dementia are sometimes treated with antipsychotics, but this has an increased risk of early death.
As of 2020, there were approximately 50 million people worldwide with Alzheimer's disease. It most often begins in people over 65 years of age, although up to 10% of cases are early-onset impacting those in their 30s to mid-60s. It affects about 6% of people 65 years and older, and women more often than men. The disease is named after German psychiatrist and pathologist Alois Alzheimer, who first described it in 1906. Alzheimer's financial burden on society is large, with an estimated global annual cost of trillion. It is ranked as the seventh leading cause of death worldwide.
Given the widespread impacts of Alzheimer's disease, both basic-science and health funders in many countries support Alzheimer's research at large scales. For example, the US National Institutes of Health program for Alzheimer's research, the National Plan to Address Alzheimer’s Disease, has a budget of US$3.98 billion for fiscal year 2026. In the European Union, the 2020 Horizon Europe research programme awarded over €570 million for dementia-related projects.
Signs and symptoms
The course of Alzheimer's is generally described in three stages, with a progressive pattern of cognitive and functional impairment. The three stages are described as early or mild, middle or moderate, and late or severe. The disease is known to target the hippocampus which is associated with memory, and this is responsible for the first symptoms of memory impairment. As the disease progresses so does the degree of memory impairment.
First symptoms
The first symptoms are often mistakenly attributed to aging or stress. Detailed neuropsychological testing can reveal mild cognitive difficulties up to eight years before a person fulfills the clinical criteria for diagnosis of Alzheimer's disease. These early symptoms can affect the most complex activities of daily living. The most noticeable deficit is short term memory loss, which shows up as difficulty in remembering recently learned facts and inability to acquire new information.
Subtle problems with the executive functions of attentiveness, planning, flexibility, and abstract thinking, or impairments in semantic memory (memory of meanings, and concept relationships) can also be symptomatic of the early stages of Alzheimer's disease. Apathy and depression can be seen at this stage, with apathy remaining as the most persistent symptom throughout the course of the disease. People with objective signs of cognitive impairment, but not more severe symptoms, may be diagnosed with mild cognitive impairment (MCI). If memory loss is the predominant symptom of MCI, it is termed amnestic MCI and is frequently seen as a prodromal or early stage of Alzheimer's disease. Amnestic MCI has a greater than 90% likelihood of being associated with Alzheimer's.
Early stage
In people with Alzheimer's disease, the increasing impairment of learning and memory eventually leads to a definitive diagnosis. In a small percentage, difficulties with language, executive functions, perception (agnosia), or execution of movements (apraxia) are more prominent than memory problems. Alzheimer's disease does not affect all memory capacities equally. Older memories of the person's life (episodic memory), facts learned (semantic memory), and implicit memory (the memory of the body on how to do things, such as using a fork to eat or how to drink from a glass) are affected to a lesser degree than new facts or memories.
Language problems are mainly characterised by a shrinking vocabulary and decreased word fluency, leading to a general impoverishment of oral and written language. In this stage, the person with Alzheimer's is usually capable of communicating basic ideas adequately. While performing fine motor tasks such as writing, drawing, or dressing, certain movement coordination and planning difficulties (apraxia) may be present; however, they are commonly unnoticed. As the disease progresses, people with Alzheimer's disease can often continue to perform many tasks independently; however, they may need assistance or supervision with the most cognitively demanding activities.
Middle stage
Progressive deterioration eventually hinders independence, with subjects being unable to perform most common activities of daily living. Speech difficulties become evident due to an inability to recall vocabulary, which leads to frequent incorrect word substitutions (paraphasias). Reading and writing skills are also progressively lost. Complex motor sequences become less coordinated as time passes and Alzheimer's disease progresses, so the risk of falling increases. During this phase, memory problems worsen, and the person may fail to recognise close relatives. Long-term memory, which was previously intact, becomes impaired.
Behavioral and neuropsychiatric changes become more prevalent. Common manifestations are wandering, irritability and emotional lability, leading to crying, outbursts of unpremeditated aggression, or resistance to caregiving. Sundowning can also appear. Approximately 30% of people with Alzheimer's disease develop illusionary misidentifications and other delusional symptoms. Subjects also lose insight of their disease process and limitations (anosognosia). Urinary incontinence can develop. These symptoms create stress for relatives and caregivers, which can be reduced by moving the person from home care to other long-term care facilities.
Late stage
During the final stage, known as the late-stage or severe stage, there is complete dependence on caregivers. Language is reduced to simple phrases or even single words, eventually leading to complete loss of speech. Despite the loss of verbal language abilities, people can often understand and return emotional signals. Although aggressiveness can still be present, extreme apathy and exhaustion are much more common symptoms. People with Alzheimer's disease will ultimately not be able to perform even the simplest tasks independently; muscle mass and mobility deteriorates to the point where they are bedridden and unable to feed themselves. The cause of death is usually an external factor, such as infection of pressure ulcers or pneumonia, not the disease itself. In some cases, there is a paradoxical lucidity immediately before death, where there is an unexpected recovery of mental clarity.
Causes
Alzheimer's disease is believed to occur when abnormal amounts of amyloid beta (Aβ), accumulating extracellularly as amyloid plaques and tau proteins, or intracellularly as neurofibrillary tangles, form in the brain, affecting neuronal functioning and connectivity, resulting in a progressive loss of brain function. This altered protein clearance ability is age-related, regulated by brain cholesterol, and associated with other neurodegenerative diseases.
The cause for most Alzheimer's cases is still mostly unknown, except for 1–2% of cases where deterministic genetic differences have been identified. Several competing hypotheses attempt to explain the underlying cause; the most predominant hypothesis is the amyloid beta (Aβ) hypothesis.
The oldest hypothesis, on which most drug therapies are based, is the cholinergic hypothesis, which proposes that Alzheimer's disease is caused by reduced synthesis of the neurotransmitter acetylcholine. The loss of cholinergic neurons noted in the limbic system and cerebral cortex, is a key feature in the progression of Alzheimer's. The 1991 amyloid hypothesis postulated that extracellular amyloid beta (Aβ) deposits are the fundamental cause of the disease. Support for this postulate comes from the location of the gene for the amyloid precursor protein (APP) on chromosome 21, together with the fact that people with trisomy 21 (Down syndrome) who have an extra gene copy almost universally exhibit at least the earliest symptoms of Alzheimer's disease by 40 years of age. A specific isoform of apolipoprotein, APOE4, is a major genetic risk factor for Alzheimer's disease. While apolipoproteins enhance the breakdown of beta amyloid, some isoforms are not very effective at this task (such as APOE4), leading to excess amyloid buildup in the brain.
Genetic
Late onset
Late-onset Alzheimer's is about 70% heritable. Genetic models in 2020 predict Alzheimer's disease with 90% accuracy. Most cases of Alzheimer's are not familial, and so they are termed sporadic Alzheimer's disease. Of the cases of sporadic Alzheimer's disease, most are classified as late onset where they are developed after the age of 65 years.
The strongest genetic risk factor for sporadic Alzheimer's disease is APOEε4. APOEε4 is one of four alleles of apolipoprotein E (APOE). APOE plays a major role in lipid-binding proteins in lipoprotein particles and the ε4 allele disrupts this function. Between 40% and 80% of people with Alzheimer's disease possess at least one APOEε4 allele. The APOEε4 allele increases the risk of the disease by three times in heterozygotes and by 15 times in homozygotes. Like many human diseases, environmental effects and genetic modifiers result in incomplete penetrance. For example, Nigerian Yoruba people do not show the relationship between dose of APOEε4 and incidence or age-of-onset for Alzheimer's disease seen in other human populations.
Early onset
Only 1–2% of Alzheimer's cases are inherited due to autosomal dominant effects, as Alzheimer's is highly polygenic. When the disease is caused by autosomal dominant variants, it is known as early onset familial Alzheimer's disease, which is rarer and has a faster rate of progression. Less than 5% of sporadic Alzheimer's disease have an earlier onset, and early-onset Alzheimer's is about 90% heritable. Familial Alzheimer's disease usually implies two or more persons affected in one or more generations.
Early onset familial Alzheimer's disease can be attributed to mutations in one of three genes: those encoding amyloid-beta precursor protein (APP) and presenilins PSEN1 and PSEN2. Most mutations in the APP and presenilin genes increase the production of a small protein called amyloid beta (Aβ)42, which is the main component of amyloid plaques. Some of the mutations merely alter the ratio between Aβ42 and the other major forms—particularly Aβ40—without increasing Aβ42 levels in the brain. Two other genes associated with autosomal dominant Alzheimer's disease are ABCA7 and SORL1.
Alleles in the TREM2 gene have been associated with a three to five times higher risk of developing Alzheimer's disease.
A Japanese pedigree of familial Alzheimer's disease was found to be associated with a deletion mutation of codon 693 of APP. This mutation and its association with Alzheimer's disease was first reported in 2008, and is known as the Osaka mutation. Only homozygotes with this mutation have an increased risk of developing Alzheimer's disease. This mutation accelerates Aβ oligomerization but the proteins do not form the amyloid fibrils that aggregate into amyloid plaques, suggesting that it is the Aβ oligomerization rather than the fibrils that may be the cause of this disease. Mice expressing this mutation have all the usual pathologies of Alzheimer's disease.
Hypotheses
Amyloid beta and tau protein
The tau hypothesis proposes that tau protein abnormalities initiate the disease cascade. In this model, hyperphosphorylated tau begins to pair with other threads of tau as paired helical filaments. Eventually, they form neurofibrillary tangles inside neurons. When this occurs, the microtubules disintegrate, destroying the structure of the cell's cytoskeleton which collapses the neuron's transport system.
A number of studies connect the misfolded amyloid beta and tau proteins associated with the pathology of Alzheimer's disease, as bringing about oxidative stress that leads to neuroinflammation. This chronic inflammation is also a feature of other neurodegenerative diseases including Parkinson's disease, and ALS. Spirochete infections have also been linked to dementia. DNA damages accumulate in Alzheimer's diseased brains; reactive oxygen species may be the major source of this DNA damage.
Sleep
Sleep disturbances are seen as a possible risk factor for inflammation in Alzheimer's disease. Sleep disruption was previously only seen as a consequence of Alzheimer's disease, but , accumulating evidence suggests that this relationship may be bidirectional.
Metal toxicity, smoking, neuroinflammation and air pollution
The cellular homeostasis of biometals such as ionic copper, iron, and zinc is disrupted in Alzheimer's disease, though it remains unclear whether this is produced by or causes the changes in proteins. Smoking is a significant Alzheimer's disease risk factor. Systemic markers of the innate immune system are risk factors for late-onset Alzheimer's disease. Exposure to air pollution may be a contributing factor to the development of Alzheimer's disease.
Age-related myelin decline
Retrogenesis is a medical hypothesis that just as the fetus goes through a process of neurodevelopment beginning with neurulation and ending with myelination, the brains of people with Alzheimer's disease go through a reverse neurodegeneration process starting with demyelination and death of axons (white matter) and ending with the death of grey matter. Likewise the hypothesis is, that as infants go through states of cognitive development, people with Alzheimer's disease go through the reverse process of progressive cognitive impairment.
According to one theory, dysfunction of oligodendrocytes and their associated myelin during aging contributes to axon damage, which in turn generates in amyloid production and tau hyperphosphorylation. Comorbidities between the demyelinating disease, multiple sclerosis, and Alzheimer's disease have been reported.
Other hypotheses
The association with celiac disease is unclear, with a 2019 study finding no increase in dementia overall in those with celiac disease while a 2018 review found an association with several types of dementia including Alzheimer's disease.
Studies have shown a potential link between infection with certain viruses and developing Alzheimer's disease later in life. Notably, a large scale study conducted on 6,245,282 patients has shown an increased risk of developing Alzheimer's disease following COVID-19 infection in cognitively normal individuals over 65.
Some evidence suggests that some viral infections such as Herpes simplex virus 1 (HSV-1) may be associated with dementia, but there are conflicting results and the association with Alzheimer's is unclear as of 2024.
Pathophysiology
Neuropathology
Alzheimer's disease is characterised by loss of neurons and synapses in the cerebral cortex and certain subcortical regions. This loss results in gross atrophy of the affected regions, including degeneration in the temporal lobe and parietal lobe, and parts of the frontal cortex and cingulate gyrus. Degeneration is also present in brainstem nuclei particularly the locus coeruleus in the pons. Studies using MRI and PET have documented reductions in the size of specific brain regions in people with Alzheimer's disease as they progressed from mild cognitive impairment to Alzheimer's disease, and in comparison with similar images from healthy older adults.
Both Aβ plaques and neurofibrillary tangles are clearly visible by microscopy in brains of those with Alzheimer's disease, especially in the hippocampus. However, Alzheimer's disease may occur without neurofibrillary tangles in the neocortex. Plaques are dense, mostly insoluble deposits of beta-amyloid peptide and cellular material outside and around neurons. Neurofibrillary tangles are aggregates of the microtubule-associated protein tau which has become hyperphosphorylated and accumulate inside the cells themselves. Although many older individuals develop some plaques and tangles as a consequence of aging, the brains of people with Alzheimer's disease have a greater number of them in specific brain regions such as the temporal lobe. Lewy bodies are not rare in the brains of people with Alzheimer's disease.
Biochemistry
Amyloid beta
Alzheimer's disease has been identified as a protein misfolding disease, a proteopathy, caused by the accumulation of abnormally folded amyloid beta protein into amyloid plaques, and tau protein into neurofibrillary tangles in the brain. Plaques are made up of small peptides, 39–43 amino acids in length, called amyloid beta. Amyloid beta is a fragment from the larger amyloid-beta precursor protein (APP) a transmembrane protein that penetrates the cell's membrane. APP is critical to neuron growth, survival, and post-injury repair. In Alzheimer's disease, gamma secretase and beta secretase act together in a proteolytic process which causes APP to be divided into smaller fragments. Although commonly researched as neuronal proteins, APP and its processing enzymes are abundantly expressed by other brain cells. One of these fragments gives rise to fibrils of amyloid beta, which then form clumps that deposit outside neurons in dense formations known as amyloid plaques. Excitatory neurons are known to be the major producers of amyloid beta that contribute to major extracellular plaque deposition.
Phosphorylated tau
Alzheimer's disease is also considered a tauopathy due to abnormal aggregation of the tau protein. Every neuron has a cytoskeleton, an internal support structure partly made up of structures called microtubules. These microtubules act like tracks, guiding nutrients and molecules from the body of the cell to the ends of the axon and back. A protein called tau stabilises the microtubules when phosphorylated, and is therefore called a microtubule-associated protein. In Alzheimer's disease, tau undergoes chemical changes, becoming hyperphosphorylated; it then begins to pair with other threads, creating neurofibrillary tangles and disintegrating the neuron's transport system. Pathogenic tau can also cause neuronal death through transposable element dysregulation. Necroptosis has also been reported as a mechanism of cell death in brain cells affected with tau tangles.
Disease mechanism
Exactly how disturbances of production and aggregation of the beta-amyloid peptide give rise to the pathology of Alzheimer's disease is not known. The amyloid hypothesis traditionally points to the accumulation of beta-amyloid peptides as the central event triggering neuron degeneration. Accumulation of aggregated amyloid fibrils, which are believed to be the toxic form of the protein responsible for disrupting the cell's calcium ion homeostasis, induces programmed cell death (apoptosis). It is also known that Aβ selectively builds up in the mitochondria in the cells of Alzheimer's-affected brains, and it also inhibits certain enzyme functions and the utilisation of glucose by neurons.
Evidence supports Aβ as playing a central role in the pathogenesis of AD, but it does not completely explain the condition, as individuals may have normal cognition and very high Aβ burden in their brains at an advanced age, and the beneficial effect of therapeutics (such as monoclonal antibodies) promoting Aβ clearance has ranged from nonexistent to modest.
Iron dyshomeostasis is linked to disease progression, an iron-dependent form of regulated cell death called ferroptosis could be involved. Products of lipid peroxidation are also elevated in AD brain compared with controls.
Various inflammatory processes and cytokines may also have a role in the pathology of Alzheimer's disease. Inflammation is a general marker of tissue damage in any disease, and may be either secondary to tissue damage in Alzheimer's disease or a marker of an immunological response. There is increasing evidence of a strong interaction between the neurons and the immunological mechanisms in the brain. Obesity and systemic inflammation may interfere with immunological processes which promote disease progression.
Alterations in the distribution of different neurotrophic factors and in the expression of their receptors such as the brain-derived neurotrophic factor (BDNF) have been described in Alzheimer's disease.
Evidence has accrued for microglia as central actors in the mechanism of AD. Microglia are topographically associated with pTau and Aβ within the brain, even when each pathologic component occurs in distinct brain regions, and microglial activation has been documented in those with mild cognitive impairment, despite a lack of tracer uptake, suggesting that microglial dysfunction may precede plaque deposition as an inciting event in AD. Microglia are the principal immunological cells of the central nervous system, serving as the tissue-resident macrophages of the brain; they are capable of recognizing and taking up Aβ through multiple pattern recognition receptors, making them central to amyloid clearance within the brain. However, microglia can also be a major source of pro-inflammatory mediators which can be deleterious to neurological function.
Diagnosis
Alzheimer's disease (AD) can only be definitively diagnosed with autopsy findings; in the absence of autopsy, clinical diagnoses of AD are "possible" or "probable", based on other findings. Up to 23% of those clinically diagnosed with AD may be misdiagnosed and may have pathology suggestive of another condition with symptoms that mimic those of AD.
AD is usually clinically diagnosed based on a person's medical history, observations from friends or relatives, and behavioral changes. The presence of characteristic neuropsychological changes with impairments in at least two cognitive domains that are severe enough to affect a person's functional abilities are required for the diagnosis. Domains that may be impaired include memory (most commonly impaired), language, executive function, visuospatial functioning, or other areas of cognition. The neurocognitive changes must be a decline from a prior level of function and the diagnosis requires ruling out other common causes of neurocognitive decline. Advanced medical imaging with computed tomography (CT) or magnetic resonance imaging (MRI), and with single-photon emission computed tomography (SPECT) or positron emission tomography (PET), can be used to help exclude other cerebral pathology or subtypes of dementia. On MRI or CT, Alzheimer's disease usually shows a generalized or focal cortical atrophy, which may be asymmetric. Atrophy of the hippocampus is also commonly seen. Brain imaging commonly also shows cerebrovascular disease, most commonly previous strokes (small or large territory strokes), and this is thought to be a contributing cause of many cases of dementia (up to 46% cases of dementia also have cerebrovascular disease on imaging). FDG-PET scan is not required for the diagnosis but it is sometimes used when standard testing is unclear. FDG-PET shows a bilateral, asymetric, temporal and parietal reduced activity. Advanced imaging may predict conversion from prodromal stages (mild cognitive impairment) to Alzheimer's disease. FDA-approved radiopharmaceutical diagnostic agents used in PET for Alzheimer's disease are florbetapir (2012), flutemetamol (2013), florbetaben (2014), and flortaucipir (2020). Because many insurance companies in the United States do not cover this procedure, its use in clinical practice is largely limited to clinical trials .
Assessment of intellectual functioning including memory testing can further characterise the state of the disease. Medical organizations have created diagnostic criteria to ease and standardise the diagnostic process for practising physicians. Definitive diagnosis can only be confirmed with post-mortem evaluations when brain material is available and can be examined histologically for senile plaques and neurofibrillary tangles.
Criteria
There are three sets of criteria for the clinical diagnoses of the spectrum of Alzheimer's disease: the 2013 fifth edition of the Diagnostic and Statistical Manual of Mental Disorders (DSM-5); the National Institute on Aging-Alzheimer's Association (NIA-AA) definition as revised in 2011; and the International Working Group criteria as revised in 2010.
Eight intellectual domains are most commonly impaired in AD—memory, language, perceptual skills, attention, motor skills, orientation, problem solving and executive functional abilities, as listed in the fourth text revision of the DSM (DSM-IV-TR).
The DSM-5 defines criteria for probable or possible AD for both major and mild neurocognitive disorder. Major or mild neurocognitive disorder must be present along with at least one cognitive deficit for a diagnosis of either probable or possible AD. For major neurocognitive disorder due to AD, probable Alzheimer's disease can be diagnosed if the individual has genetic evidence of AD or if two or more acquired cognitive deficits, and a functional disability that is not from another disorder, are present. Otherwise, possible AD can be diagnosed as the diagnosis follows an atypical route. For mild neurocognitive disorder due to AD, probable Alzheimer's disease can be diagnosed if there is genetic evidence, whereas possible AD can be met if all of the following are present: no genetic evidence, decline in both learning and memory, two or more cognitive deficits, and a functional disability not from another disorder.
The NIA-AA criteria are used mainly in research rather than in clinical assessments. They define AD through three major stages: preclinical, mild cognitive impairment (MCI), and Alzheimer's dementia. Diagnosis in the preclinical stage is complex and focuses on asymptomatic individuals; the latter two stages describe individuals experiencing symptoms, along with biomarkers, predominantly those for neuronal injury (mainly tau-related) and amyloid beta deposition. The core clinical criteria itself rests on the presence of cognitive impairment without the presence of comorbidities. The third stage is divided into probable and possible AD dementia. In probable AD dementia there is steady impairment of cognition over time and a memory-related or non-memory-related cognitive dysfunction. In possible AD dementia, another causal disease such as cerebrovascular disease is present.
Techniques
Neuropsychological tests including cognitive tests such as the mini–mental state examination (MMSE), the Montreal Cognitive Assessment (MoCA) and the Mini-Cog are widely used to aid in diagnosis of the cognitive impairments in AD. These tests may not always be accurate, as they lack sensitivity to mild cognitive impairment, and can be biased by language or attention problems; more comprehensive test arrays are necessary for high reliability of results, particularly in the earliest stages of the disease.
Further neurological examinations are crucial in the differential diagnosis of Alzheimer's disease and other diseases. Interviews with family members are used in assessment; caregivers can supply important information on daily living abilities and on the decrease in the person's mental function. A caregiver's viewpoint is particularly important, since a person with Alzheimer's disease is commonly unaware of their deficits. Many times, families have difficulties in the detection of initial dementia symptoms and may not communicate accurate information to a physician.
Supplemental testing can rule out other potentially treatable diagnoses and help avoid misdiagnoses. Common supplemental tests include blood tests, thyroid function tests, as well as tests to assess vitamin B12 levels, rule out neurosyphilis and rule out metabolic problems (including tests for kidney function, electrolyte levels and for diabetes). MRI or CT scans might also be used to rule out other potential causes of the symptoms – including tumors or strokes. Delirium and depression can be common among individuals and are important to rule out.
Psychological tests for depression are used, since depression can either be concurrent with AD (see Depression of Alzheimer disease), an early sign of cognitive impairment, or even the cause.
Due to low accuracy, the C-PIB-PET scan is not recommended as an early diagnostic tool or for predicting the development of AD when people show signs of mild cognitive impairment (MCI). The use of 18F-FDG PET scans, as a single test, to identify people who may develop Alzheimer's disease is not supported by evidence.
Prevention
There are no disease-modifying treatments available to cure Alzheimer's disease and because of this, AD research has focused on interventions to prevent the onset and progression. There is no evidence that supports any particular measure in preventing AD, and studies of measures to prevent the onset or progression have produced inconsistent results. Epidemiological studies have proposed relationships between an individual's likelihood of developing AD and modifiable factors, such as medications, lifestyle, and diet. There are some challenges in determining whether interventions for AD act as a primary prevention method, preventing the disease itself, or a secondary prevention method, identifying the early stages of the disease. These challenges include duration of intervention, different stages of disease at which intervention begins, and lack of standardization of inclusion criteria regarding biomarkers specific for AD. Further research is needed to determine factors that can help prevent AD.
Medication
Cardiovascular risk factors, such as hypercholesterolaemia, hypertension, diabetes, and smoking, are associated with a higher risk of onset and worsened course of AD. The use of statins to lower cholesterol may be of benefit in AD. Antihypertensive and antidiabetic medications in individuals without overt cognitive impairment may decrease the risk of dementia by influencing cerebrovascular pathology. More research is needed to examine the relationship with AD specifically; clarification of the direct role medications play versus other concurrent lifestyle changes (diet, exercise, smoking) is needed.
Depression is associated with an increased risk for AD; management with antidepressant medications may provide a preventative measure.
Historically, long-term usage of non-steroidal anti-inflammatory drugs (NSAIDs) were thought to be associated with a reduced likelihood of developing AD as it reduces inflammation, but NSAIDs do not appear to be useful as a treatment. Additionally, because women have a higher incidence of AD than men, it was once thought that estrogen deficiency during menopause was a risk factor, but there is a lack of evidence to show that hormone replacement therapy (HRT) in menopause decreases risk of cognitive decline.
Lifestyle
Certain lifestyle activities, such as physical and cognitive exercises, higher education and occupational attainment, cigarette smoking, stress, sleep, and the management of other comorbidities, including diabetes and hypertension, may affect the risk of developing AD.
Physical exercise is associated with a decreased rate of dementia, and is effective in reducing symptom severity in those with AD. Memory and cognitive functions can be improved with aerobic exercises including brisk walking three times weekly for forty minutes. It may also induce neuroplasticity of the brain. Participating in mental exercises, such as reading, crossword puzzles, and chess have shown potential to be preventive. Meeting the WHO recommendations for physical activity is associated with a lower risk of AD.
Higher education and occupational attainment, and participation in leisure activities, contribute to a reduced risk of developing AD, or of delaying the onset of symptoms. This is compatible with the cognitive reserve theory, which states that some life experiences result in more efficient neural functioning providing the individual a cognitive reserve that delays the onset of dementia manifestations. Education delays the onset of Alzheimer's disease syndrome without changing the duration of the disease.
Cessation in smoking may reduce risk of developing AD, specifically in those who carry the APOE ɛ4 allele. The increased oxidative stress caused by smoking results in downstream inflammatory or neurodegenerative processes that may increase risk of developing AD. Avoidance of smoking, counseling and pharmacotherapies to quit smoking are used, and avoidance of environmental tobacco smoke is recommended.
Alzheimer's disease is associated with sleep disorders but the precise relationship is unclear. It was once thought that as people get older, the risk of developing sleep disorders and AD independently increase, but research suggests sleep disorders may be a risk factor for AD. One theory is that the mechanisms to increase clearance of toxic substances, including Aβ, are active during sleep. With decreased sleep, a person is increasing Aβ production and decreasing Aβ clearance, resulting in Aβ accumulation. Receiving adequate sleep (approximately 7–8 hours) every night has become a potential lifestyle intervention to prevent the development of AD.
Stress is a risk factor for the development of AD. The mechanism by which stress predisposes someone to development of AD is unclear, but it is suggested that lifetime stressors may affect a person's epigenome, leading to an overexpression or under expression of specific genes. Although the relationship of stress and AD is unclear, strategies to reduce stress and relax the mind may be helpful strategies in preventing the progression or Alzheimer's disease. Meditation, for instance, is a helpful lifestyle change to support cognition and well-being, though further research is needed to assess long-term effects.
Management
There is no cure for AD; available treatments offer relatively small symptomatic benefits but remain palliative in nature. Treatments can be divided into pharmaceutical, psychosocial, and caregiving.
Pharmaceutical
Medications used to treat the cognitive symptoms of AD rather than the underlying cause include: four acetylcholinesterase inhibitors (tacrine, rivastigmine, galantamine, and donepezil) and memantine, an NMDA receptor antagonist. The acetylcholinesterase inhibitors are intended for those with mild to severe AD, whereas memantine is intended for those with moderate or severe Alzheimer's disease. The benefit from their use is small.
Reduction in the activity of the cholinergic neurons is a well-known feature of AD. Acetylcholinesterase inhibitors are employed to reduce the rate at which acetylcholine (ACh) is broken down, thereby increasing the concentration of ACh in the brain and combating the loss of ACh caused by the death of cholinergic neurons. There is evidence for the efficacy of these medications in mild to moderate AD, and some evidence for their use in the advanced stage. The use of these drugs in mild cognitive impairment has not shown any effect in a delay of the onset of Alzheimer's disease. The most common side effects are nausea and vomiting, both of which are linked to cholinergic excess. These side effects arise in approximately 10–20% of users, are mild to moderate in severity, and can be managed by slowly adjusting medication doses. Less common secondary effects include muscle cramps, decreased heart rate (bradycardia), decreased appetite and weight, and increased gastric acid production.
Glutamate is an excitatory neurotransmitter of the nervous system, although excessive amounts in the brain can lead to cell death through a process called excitotoxicity which consists of the overstimulation of glutamate receptors. Excitotoxicity occurs not only in AD, but also in other neurological diseases such as Parkinson's disease and multiple sclerosis. Memantine is a noncompetitive NMDA receptor antagonist first used as an anti-influenza agent. It acts on the glutamatergic system by blocking NMDA receptors and inhibiting their overstimulation by glutamate. Memantine has been shown to have a small benefit in the treatment of moderate to severe AD. Reported adverse events with memantine are infrequent and mild, including hallucinations, confusion, dizziness, headache and fatigue. The combination of memantine and donepezil has been shown to be "of statistically significant but clinically marginal effectiveness".
An extract of Ginkgo biloba known as EGb 761 has been used for treating AD and other neuropsychiatric disorders. Its use is approved throughout Europe. The World Federation of Biological Psychiatry guidelines lists EGb 761 with the same weight of evidence (level B) given to acetylcholinesterase inhibitors and memantine. EGb 761 is the only one that showed improvement of symptoms in both AD and vascular dementia. EGb 761 may have a role either on its own or as an add-on if other therapies prove ineffective. A 2016 review concluded that the quality of evidence from clinical trials on Ginkgo biloba has been insufficient to warrant its use for treating AD.
Atypical antipsychotics are modestly useful in reducing aggression and psychosis in people with AD, but their advantages are offset by serious adverse effects, such as stroke, movement difficulties or cognitive decline. When used in the long-term, they have been shown to associate with increased mortality. They are recommended in dementia only after first line therapies such as behavior modification have failed, and due to the risk of adverse effects, they should be used for the shortest amount of time possible. Stopping antipsychotic use in this group of people appears to be safe.
Psychosocial
Psychosocial interventions are used as an adjunct to pharmaceutical treatment and can be classified within behavior-, emotion-, cognition- or stimulation-oriented approaches.
Behavioral interventions attempt to identify and reduce the antecedents and consequences of problem behaviors. This approach has not shown success in improving overall functioning, but can help to reduce some specific problem behaviors, such as incontinence. There is a lack of high quality data on the effectiveness of these techniques in other behavior problems such as wandering. Music therapy is effective in reducing behavioral and psychological symptoms.
Emotion-oriented interventions include reminiscence therapy, validation therapy, supportive psychotherapy, sensory integration, also called snoezelen, and simulated presence therapy. A Cochrane review has found no evidence that this is effective. Reminiscence therapy (RT) involves the discussion of past experiences individually or in group, many times with the aid of photographs, household items, music and sound recordings, or other familiar items from the past. A 2018 review of the effectiveness of RT found that effects were inconsistent, small in size and of doubtful clinical significance, and varied by setting. Simulated presence therapy (SPT) is based on attachment theories and involves playing a recording with voices of the closest relatives of the person with AD. There is partial evidence indicating that SPT may reduce challenging behaviors.
The aim of cognition-oriented treatments, which include reality orientation and cognitive retraining, is the reduction of cognitive deficits. Reality orientation consists of the presentation of information about time, place, or person to ease the understanding of the person about its surroundings and his or her place in them. On the other hand, cognitive retraining tries to improve impaired capacities by exercising mental abilities. Both have shown some efficacy improving cognitive capacities.
Stimulation-oriented treatments include art, music and pet therapies, exercise, and any other kind of recreational activities. Stimulation has modest support for improving behavior, mood, and, to a lesser extent, function. Nevertheless, as important as these effects are, the main support for the use of stimulation therapies is the change in the person's routine.
Caregiving
Since AD has no cure and it gradually renders people incapable of tending to their own needs, caregiving is essentially the treatment and must be carefully managed over the course of the disease.
During the early and moderate stages, modifications to the living environment and lifestyle can increase safety and reduce caretaker burden. Examples of such modifications are the adherence to simplified routines, the placing of safety locks, the labeling of household items to cue the person with the disease or the use of modified daily life objects. If eating becomes problematic, food will need to be prepared in smaller pieces or even puréed. When swallowing difficulties arise, the use of feeding tubes may be required. In such cases, the medical efficacy and ethics of continuing feeding is an important consideration of the caregivers and family members. The use of physical restraints is rarely indicated in any stage of the disease, although there are situations when they are necessary to prevent harm to the person with Alzheimer's disease or their caregivers.
During the final stages of the disease, treatment is centred on relieving discomfort until death, often with the help of hospice.
Diet
Diet may be a modifiable risk factor for the development of Alzheimer's disease but more research needs to be conducted. The Mediterranean diet, and the DASH diet are both associated with less cognitive decline. A different approach has been to incorporate elements of both of these diets into one known as the MIND diet.
Results from large-scale epidemiological studies and clinical trials have not demonstrated an independent role for most individual dietary components. Low folate has been associated with an increased risk of Alzheimer's disease.
Prognosis
The early stages of AD are difficult to diagnose. A definitive diagnosis is usually made once cognitive impairment compromises daily living activities, although the person may still be living independently. The symptoms will progress from mild cognitive problems, such as memory loss through increasing stages of cognitive and non-cognitive disturbances, eliminating any possibility of independent living, especially in the late stages of the disease.
Life expectancy of people with AD is reduced. The normal life expectancy for 60 to 70 years old is 23 to 15 years; for 90 years old it is 4.5 years. Following AD diagnosis it ranges from 7 to 10 years for those in their 60s and early 70s (a loss of 13 to 8 years), to only about 3 years or less (a loss of 1.5 years) for those in their 90s.
Fewer than 3% of people live more than fourteen years after diagnosis. Disease features significantly associated with reduced survival are an increased severity of cognitive impairment, decreased functional level, disturbances in the neurological examination, history of falls, malnutrition, dehydration and weight loss. Other coincident diseases such as heart problems, diabetes, or history of alcohol abuse are also related with shortened survival. While the earlier the age at onset the higher the total survival years, life expectancy is particularly reduced when compared to the healthy population among those who are younger. Men have a less favourable survival prognosis than women.
Aspiration pneumonia is the most frequent immediate cause of death brought by AD. While the reasons behind the lower prevalence of cancer in AD patients remain unclear, some researchers hypothesize that biological mechanisms shared by both diseases might play a role. However, this requires further investigation.
Epidemiology
Two main measures are used in epidemiological studies: incidence and prevalence. Incidence is the number of new cases per unit of person-time at risk (usually number of new cases per thousand person-years); while prevalence is the total number of cases of the disease in the population at any given time.
Regarding incidence, cohort longitudinal studies (studies where a disease-free population is followed over the years) provide rates between 10 and 15 per thousand person-years for all dementias and 5–8 for AD, which means that half of new dementia cases each year are Alzheimer's disease. Advancing age is a primary risk factor for the disease and incidence rates are not equal for all ages: every 5 years after the age of 65, the risk of acquiring the disease approximately doubles, increasing from 3 to as much as 69 per thousand person years. Females with AD are more common than males, but this difference is likely due to women's longer life spans. When adjusted for age, both sexes are affected by Alzheimer's at equal rates. In the United States, the risk of dying from AD in 2010 was 26% higher among the non-Hispanic white population than among the non-Hispanic black population, and the Hispanic population had a 30% lower risk than the non-Hispanic white population. However, much AD research remains to be done in minority groups, such as the African American, East Asian and Hispanic/Latino populations. Studies have shown that these groups are underrepresented in clinical trials and do not have the same risk of developing AD when carrying certain genetic risk factors (i.e. APOE4), compared to their caucasian counterparts.
The prevalence of AD in populations is dependent upon factors including incidence and survival. Since the incidence of AD increases with age, prevalence depends on the mean age of the population for which prevalence is given. In the United States in 2020, AD dementia prevalence was estimated to be 5.3% for those in the 60–74 age group, with the rate increasing to 13.8% in the 74–84 group and to 34.6% in those greater than 85. Prevalence rates in some less developed regions around the globe are lower. Both the prevalence and incidence rates of AD are steadily increasing, and the prevalence rate is estimated to triple by 2050 reaching 152 million, compared to the 50 million people with AD globally in 2020.
History
The ancient Greek and Roman philosophers and physicians associated old age with increasing dementia. It was not until 1901 that German psychiatrist Alois Alzheimer identified the first case of what became known as Alzheimer's disease, named after him, in a fifty-year-old woman he called Auguste D. He followed her case until she died in 1906 when he first reported publicly on it. During the next five years, eleven similar cases were reported in the medical literature, some of them already using the term Alzheimer's disease. The disease was first described as a distinctive disease by Emil Kraepelin after suppressing some of the clinical (delusions and hallucinations) and pathological features (arteriosclerotic changes) contained in the original report of Auguste D. He included Alzheimer's disease, also named presenile dementia by Kraepelin, as a subtype of senile dementia in the eighth edition of his Textbook of Psychiatry, published on 15 July 1910.
For most of the 20th century, the diagnosis of Alzheimer's disease was reserved for individuals between the ages of 45 and 65 who developed symptoms of dementia. The terminology changed after 1977 when a conference on Alzheimer's disease concluded that the clinical and pathological manifestations of presenile and senile dementia were almost identical, although the authors also added that this did not rule out the possibility that they had different causes. This eventually led to the diagnosis of Alzheimer's disease independent of age. The term senile dementia of the Alzheimer type (SDAT) was used for a time to describe the condition in those over 65, with classical Alzheimer's disease being used to describe those who were younger. Eventually, the term Alzheimer's disease was formally adopted in medical nomenclature to describe individuals of all ages with a characteristic common symptom pattern, disease course, and neuropathology.
The National Institute of Neurological and Communicative Disorders and Stroke (NINCDS) and the Alzheimer's Disease and Related Disorders Association (ADRDA, now known as the Alzheimer's Association) established the most commonly used NINCDS-ADRDA Alzheimer's Criteria for diagnosis in 1984, extensively updated in 2007. These criteria require that the presence of cognitive impairment, and a suspected dementia syndrome, be confirmed by neuropsychological testing for a clinical diagnosis of possible or probable Alzheimer's disease. A histopathologic confirmation including a microscopic examination of brain tissue is required for a definitive diagnosis. Good statistical reliability and validity have been shown between the diagnostic criteria and definitive histopathological confirmation.
Society and culture
Social costs
Dementia, and specifically Alzheimer's disease, may be among the most costly diseases for societies worldwide. As populations age, these costs will probably increase and become an important social problem and economic burden. Costs associated with AD include direct and indirect medical costs, which vary between countries depending on social care for a person with AD. Direct costs include doctor visits, hospital care, medical treatments, nursing home care, specialized equipment, and household expenses. Indirect costs include the cost of informal care and the loss in productivity of informal caregivers.
In the United States , informal (family) care is estimated to constitute nearly three-fourths of caregiving for people with AD at a cost of US$234 billion per year and approximately 18.5 billion hours of care. The cost to society worldwide to care for individuals with AD is projected to increase nearly ten-fold, and reach about US$9.1 trillion by 2050.
Costs for those with more severe dementia or behavioral disturbances are higher and are related to the additional caregiving time to provide physical care.
Caregiving burden
Individuals with Alzheimer's will require assistance in their lifetime, and care will most likely come in the form of a full-time caregiver which is often a role that is taken on by the spouse or a close relative. Caregiving tends to include physical and emotional burdens as well as time and financial strain at times on the person administering the aid. Alzheimer's disease is known for placing a great burden on caregivers which includes social, psychological, physical, or economic aspects. Home care is usually preferred by both those people with Alzheimer's disease as well as their families. This option also delays or eliminates the need for more professional and costly levels of care. Nevertheless, two-thirds of nursing home residents have dementias.
Dementia caregivers are subject to high rates of physical and mental disorders. Factors associated with greater psychosocial problems of the primary caregivers include having an affected person at home, the caregiver being a spouse, demanding behaviors of the cared person such as depression, behavioral disturbances, hallucinations, sleep problems or walking disruptions and social isolation. In the United States, the yearly cost of caring for a person with dementia ranges from $28,078-$56,022 per year for formal medical care and $36,667-$92,689 for informal care provided by a relative or friend (assuming market value replacement costs for the care provided by the informal caregiver) and $15,792-$71,813 in lost wages.
Cognitive behavioral therapy and the teaching of coping strategies either individually or in group have demonstrated their efficacy in improving caregivers' psychological health.
Media
Alzheimer's disease has been portrayed in films such as: Iris (2001), based on John Bayley's memoir of his wife Iris Murdoch; The Notebook (2004), based on Nicholas Sparks's 1996 novel of the same name; A Moment to Remember (2004); Thanmathra (2005); Memories of Tomorrow (Ashita no Kioku) (2006), based on Hiroshi Ogiwara's novel of the same name; Away from Her (2006), based on Alice Munro's short story The Bear Came over the Mountain; Still Alice (2014), about a Columbia University professor who has early onset Alzheimer's disease, based on Lisa Genova's 2007 novel of the same name and featuring Julianne Moore in the title role. Documentaries on Alzheimer's disease include Malcolm and Barbara: A Love Story (1999) and Malcolm and Barbara: Love's Farewell (2007), both featuring Malcolm Pointon.
Alzheimer's disease has also been portrayed in music by English musician the Caretaker in releases such as Persistent Repetition of Phrases (2008), An Empty Bliss Beyond This World (2011), and Everywhere at the End of Time (20162019). Paintings depicting the disorder include the late works by American artist William Utermohlen, who drew self-portraits from 1995 to 2000 as an experiment of showing his disease through art.
Research directions
Antibodies may have the ability to alter the disease course by targeting amyloid beta with immunotherapy medications such as donanemab and lecanemab.
Lecanemab was approved via the FDA accelerated approval process, and was converted to traditional approval in July 2023, after further testing, along with the addition of a boxed warning about amyloid-related imaging abnormalities. As of early August 2024, lecanemab was approved for sale in Japan, South Korea, China, Hong Kong and Israel although it was recommended against approval by an advisory body of the European Union on July 26, citing its side effects.
Donanemab was approved by the FDA in July 2024. Anti-amyloid drugs also cause brain shrinkage. The cholinesterase inhibitor benzgalantamine was approved by the FDA in July 2024.
Specific medications that may reduce the risk or progression of Alzheimer's disease have been studied. The research trials investigating medications generally impact Aβ plaques, inflammation, APOE, neurotransmitter receptors, neurogenesis, growth factors or hormones.
Machine learning algorithms with electronic health records are being studied as a way to predict Alzheimer's disease earlier.
| Biology and health sciences | Non-infectious disease | null |
2740949 | https://en.wikipedia.org/wiki/Mott%20insulator | Mott insulator | Mott insulators are a class of materials that are expected to conduct electricity according to conventional band theories, but turn out to be insulators (particularly at low temperatures). These insulators fail to be correctly described by band theories of solids due to their strong electron–electron interactions, which are not considered in conventional band theory. A Mott transition is a transition from a metal to an insulator, driven by the strong interactions between electrons. One of the simplest models that can capture Mott transition is the Hubbard model.
The band gap in a Mott insulator exists between bands of like character, such as 3d electron bands, whereas the band gap in charge-transfer insulators exists between anion and cation states.
History
Although the band theory of solids had been very successful in describing various electrical properties of materials, in 1937 Jan Hendrik de Boer and Evert Johannes Willem Verwey pointed out that a variety of transition metal oxides predicted to be conductors by band theory are insulators. With an odd number of electrons per unit cell, the valence band is only partially filled, so the Fermi level lies within the band. From the band theory, this implies that such a material has to be a metal. This conclusion fails for several cases, e.g. CoO, one of the strongest insulators known.
Nevill Mott and Rudolf Peierls also in 1937 predicted the failing of band theory can be explained by including interactions between electrons.
In 1949, in particular, Mott proposed a model for NiO as an insulator, where conduction is based on the formula
(Ni2+O2−)2 → Ni3+O2− + Ni1+O2−.
In this situation, the formation of an energy gap preventing conduction can be understood as the competition between the Coulomb potential U between 3d electrons and the transfer integral t of 3d electrons between neighboring atoms (the transfer integral is a part of the tight binding approximation). The total energy gap is then
Egap = U − 2zt,
where z is the number of nearest-neighbor atoms.
In general, Mott insulators occur when the repulsive Coulomb potential U is large enough to create an energy gap. One of the simplest theories of Mott insulators is the 1963 Hubbard model. The crossover from a metal to a Mott insulator as U is increased, can be predicted within the so-called dynamical mean field theory.
Mott reviewed the subject (with a good overview) in 1968. The subject has been thoroughly reviewed in a comprehensive paper by Masatoshi Imada, Atsushi Fujimori, and Yoshinori Tokura. A recent proposal of a "Griffiths-like phase close to the Mott transition" has been reported in the literature.
Mott criterion
The Mott criterion describes the critical point of the metal–insulator transition. The criterion is
where is the electron density of the material and the effective bohr radius.
The constant , according to various estimates, is 2.0, 2.78,4.0, or 4.2.
If the criterion is satisfied (i.e. if the density of electrons is sufficiently high) the material becomes conductive (metal) and otherwise it will be an insulator.
Mottness
Mottism denotes the additional ingredient, aside from antiferromagnetic ordering, which is necessary to fully describe a Mott insulator. In other words, we might write: antiferromagnetic order + mottism = Mott insulator.
Thus, mottism accounts for all of the properties of Mott insulators that cannot be attributed simply to antiferromagnetism.
There are a number of properties of Mott insulators, derived from both experimental and theoretical observations, which cannot be attributed to antiferromagnetic ordering and thus constitute mottism. These properties include:
Spectral weight transfer on the Mott scale
Vanishing of the single particle Green function along a connected surface in momentum space in the first Brillouin zone
Two sign changes of the Hall coefficient as electron doping goes from to (band insulators have only one sign change at )
The presence of a charge (with the charge of an electron) boson at low energies
A pseudogap away from half-filling ()
Mott transition
A Mott transition is a metal-insulator transition in condensed matter. Due to electric field screening the potential energy becomes much more sharply (exponentially) peaked around the equilibrium position of the atom and electrons become localized and can no longer conduct a current. It is named after physicist Nevill Francis Mott.
Conceptual explanation
In a semiconductor at low temperatures, each 'site' (atom or group of atoms) contains a certain number of electrons and is electrically neutral. For an electron to move away from a site, it requires a certain amount of energy, as the electron is normally pulled back toward the (now positively charged) site by Coulomb forces. If the temperature is high enough that of energy is available per site, the Boltzmann distribution predicts that a significant fraction of electrons will have enough energy to escape their site, leaving an electron hole behind and becoming conduction electrons that conduct current. The result is that at low temperatures a material is insulating, and at high temperatures the material conducts.
While the conduction in an n- (p-) type doped semiconductor sets in at high temperatures because the conduction (valence) band is partially filled with electrons (holes) with the original band structure being unchanged, the situation is different in the case of the Mott transition where the band structure itself changes. Mott argued that the transition must be sudden, occurring when the density of free electrons N and the Bohr radius satisfies .
Simply put, a Mott transition is a change in a material's behavior from insulating to metallic due to various factors. This transition is known to exist in various systems: mercury metal vapor-liquid, metal NH3 solutions, transition metal chalcogenides and transition metal oxides. In the case of transition metal oxides, the material typically switches from being a good electrical insulator to a good electrical conductor. The insulator-metal transition can also be modified by changes in temperature, pressure or composition (doping). As observed by Nevill Francis Mott in his 1949 publication on Ni-oxide, the origin of this behavior is correlations between electrons and the close relationship this phenomenon has to magnetism.
The physical origin of the Mott transition is the interplay between the Coulomb repulsion of electrons and their degree of localization (band width). Once the carrier density becomes too high (e.g. due to doping), the energy of the system can be lowered by the localization of the formerly conducting electrons (band width reduction), leading to the formation of a band gap, e.g. by pressure (i.e. a semiconductor/insulator).
In a semiconductor, the doping level also affects the Mott transition. It has been observed that higher dopant concentrations in a semiconductor creates internal stresses that increase the free energy (acting as a change in pressure) of the system, thus reducing the ionization energy.
The reduced barrier causes easier transfer by tunneling or by thermal emission from donor to its adjacent donor. The effect is enhanced when pressure is applied for the reason stated previously. When the transport of carriers overcomes a minimum activation energy, the semiconductor has undergone a Mott transition and become metallic.
The Mott transition is usually first order, and involves discontinuous changes of physical properties. Theoretical studies of the Mott transition in the limit of large dimension find a first order transition. However in low dimensions and when the lattice geometry leads to frustration of magnetic ordering, it may be only weakly first order or even continuous (i.e second order). Weakly first order Mott transitions are seen in some quasi-two dimensional organic materials. Continuous Mott transitions have been reported in semiconductor moire materials. A theory of a continuous Mott transition is available if the Mott insulating phase is a quantum spin liquid with an emergent fermi surface of neutral fermions.
Applications
Mott insulators are of growing interest in advanced physics research, and are not yet fully understood. They have applications in thin-film magnetic heterostructures and the strong correlated phenomena in high-temperature superconductivity, for example.
This kind of insulator can become a conductor by changing some parameters, which may be composition, pressure, strain, voltage, or magnetic field. The effect is known as a Mott transition and can be used to build smaller field-effect transistors, switches and memory devices than possible with conventional materials.
| Physical sciences | Basics_2 | Physics |
2741037 | https://en.wikipedia.org/wiki/Hom%20functor | Hom functor | In mathematics, specifically in category theory, hom-sets (i.e. sets of morphisms between objects) give rise to important functors to the category of sets. These functors are called hom-functors and have numerous applications in category theory and other branches of mathematics.
Formal definition
Let C be a locally small category (i.e. a category for which hom-classes are actually sets and not proper classes).
For all objects A and B in C we define two functors to the category of sets as follows:
{| class=wikitable
|-
! Hom(A, –) : C → Set
! Hom(–, B) : C → Set
|-
| This is a covariant functor given by:
Hom(A, –) maps each object X in C to the set of morphisms, Hom(A, X)
Hom(A, –) maps each morphism f : X → Y to the function
Hom(A, f) : Hom(A, X) → Hom(A, Y) given by
for each g in Hom(A, X).
| This is a contravariant functor given by:
Hom(–, B) maps each object X in C to the set of morphisms, Hom(X, B)
Hom(–, B) maps each morphism h : X → Y to the function
Hom(h, B) : Hom(Y, B) → Hom(X, B) given by
for each g in Hom(Y, B).
|}
The functor Hom(–, B) is also called the functor of points of the object B.
Note that fixing the first argument of Hom naturally gives rise to a covariant functor and fixing the second argument naturally gives a contravariant functor. This is an artifact of the way in which one must compose the morphisms.
The pair of functors Hom(A, –) and Hom(–, B) are related in a natural manner. For any pair of morphisms f : B → B′ and h : A′ → A the following diagram commutes:
Both paths send g : A → B to f∘g∘h : A′ → B′.
The commutativity of the above diagram implies that Hom(–, –) is a bifunctor from C × C to Set which is contravariant in the first argument and covariant in the second. Equivalently, we may say that Hom(–, –) is a bifunctor
Hom(–, –) : Cop × C → Set
where Cop is the opposite category to C. The notation HomC(–, –) is sometimes used for Hom(–, –) in order to emphasize the category forming the domain.
Yoneda's lemma
Referring to the above commutative diagram, one observes that every morphism
h : A′ → A
gives rise to a natural transformation
Hom(h, –) : Hom(A, –) → Hom(A′, –)
and every morphism
f : B → B′
gives rise to a natural transformation
Hom(–, f) : Hom(–, B) → Hom(–, B′)
Yoneda's lemma implies that every natural transformation between Hom functors is of this form. In other words, the Hom functors give rise to a full and faithful embedding of the category C into the functor category SetCop (covariant or contravariant depending on which Hom functor is used).
Internal Hom functor
Some categories may possess a functor that behaves like a Hom functor, but takes values in the category C itself, rather than Set. Such a functor is referred to as the internal Hom functor, and is often written as
to emphasize its product-like nature, or as
to emphasize its functorial nature, or sometimes merely in lower-case:
For examples, see Category of relations.
Categories that possess an internal Hom functor are referred to as closed categories. One has that
,
where I is the unit object of the closed category. For the case of a closed monoidal category, this extends to the notion of currying, namely, that
where is a bifunctor, the internal product functor defining a monoidal category. The isomorphism is natural in both X and Z. In other words, in a closed monoidal category, the internal Hom functor is an adjoint functor to the internal product functor. The object is called the internal Hom. When is the Cartesian product , the object is called the exponential object, and is often written as .
Internal Homs, when chained together, form a language, called the internal language of the category. The most famous of these are simply typed lambda calculus, which is the internal language of Cartesian closed categories, and the linear type system, which is the internal language of closed symmetric monoidal categories.
Properties
Note that a functor of the form
Hom(–, A) : Cop → Set
is a presheaf; likewise, Hom(A, –) is a copresheaf.
A functor F : C → Set that is naturally isomorphic to Hom(A, –) for some A in C is called a representable functor (or representable copresheaf); likewise, a contravariant functor equivalent to Hom(–, A) might be called corepresentable.
Note that Hom(–, –) : Cop × C → Set is a profunctor, and, specifically, it is the identity profunctor .
The internal hom functor preserves limits; that is, sends limits to limits, while sends limits in , that is colimits in , into limits. In a certain sense, this can be taken as the definition of a limit or colimit.
The endofunctor Hom(E, –) : Set → Set can be given the structure of a monad; this monad is called the environment (or reader) monad.
Other properties
If A is an abelian category and A is an object of A, then HomA(A, –) is a covariant left-exact functor from A to the category Ab of abelian groups. It is exact if and only if A is projective.
Let R be a ring and M a left R-module. The functor HomR(M, –): Mod-R → Ab is adjoint to the tensor product functor – R M: Ab → Mod-R.
| Mathematics | Category theory | null |
2741315 | https://en.wikipedia.org/wiki/Physical%20activity | Physical activity | Physical activity is defined as any voluntary bodily movement produced by skeletal muscles that requires energy expenditure. Physical activity encompasses all activities, at any intensity, performed during any time of day or night. It includes both voluntary exercise and incidental activity integrated into the daily routine.
This integrated activity may not be planned, structured, repetitive or purposeful for the improvement of physical fitness, and may include activities such as walking to the local shop, cleaning, working, active transport etc. Lack of physical activity is associated with a range of negative health outcomes, whereas increased physical activity can improve physical and mental health, as well as cognitive and cardiovascular health. There are at least eight investments that work to increase population-level physical activity, including whole-of-school programmes, active transport, active urban design, healthcare, public education and mass media, sport for all, workplaces and community-wide programmes. Physical activity increases energy expenditure and is a key regulator in controlling body weight (see Summermatter cycle for more). In human beings, differences among individuals in the amount of physical activity have a substantial genetic basis.
Terminology misconception
"Exercise" and "physical activity" are frequently used interchangeably and generally refer to physical activity performed during leisure time with the primary purpose of improving or maintaining physical fitness, physical performance, or health. However, physical activity is not exactly the same concept as exercise. Exercise is defined as a subcategory of physical activity that is planned, structured, repetitive, and purposeful in the sense that the improvement or maintenance of one or more components of physical fitness is the objective. Conversely, physical activity includes exercise but may also be unplanned, unstructured, random and non-purposeful carried out for a multitude of reasons.
A 2021 study shows that people who start successful physical activity programmes maintain much of it for at least three months.
Intensity
Physical activity can be at any intensity, ranging from a little muscle twitch to a full-out sprint. Physical activity can be thought of as a continuum in practice, ranging from inactive lifestyles to high-intensity exercises. Intensities are broadly categorized according to energy expenditure using a standard measure of intensity, metabolic equivalents (METs). The broad categories are sedentary behavior, light activity, moderate activity, and vigorous activity.
Example activities at each intensity
The following table documents some examples of physical activities at each intensity level. Depending on the individual and the activity involved, activities may overlap intensity categories or change categories completely.
Physical activity as prevention and therapy
Physical activity is a cornerstone of public health and prevention of non-communicable disease. Physical inactivity has been found to cause a wide range of non-communicable diseases, including coronary heart disease, stroke, diabetes mellitus and depression. An analysis of the healthcare costs of non-communicable diseases and mental illness attributable to physical inactivity for 2020–30 found that 500 million new cases of disease will occur globally between 2020 and 2030 if physical activity remains at today's levels. This corresponds to more than US$300 billion in treatment costs
As of 2024, 31% of adults and 80% of adolescents do not meet the recommended levels of physical activity.
Many studies have demonstrated the potential beneficial effects of physical activity on the prevention and therapy of many disorders such as Obesity and Irritable bowel syndrome
Physical activity has been shown to reduce anxiety as a condition (individual physical exercise, without continuity), anxiety as a personality trait (continuous performance, "exercise" of certain physical activities), psycho-physiological signs of anxiety - blood pressure and heart rate (moderate physical activity can lead to a decrease in the intensity of short-term physiological reactivity and encourage recovery from short-term physiological stressors (Biddle et al., 2000)).
For people with a severe depressive episode and anxiety disorder, long and short walks proved to be the most effective; for people with substance abuse disorders, bipolar disorder and frequent psychotic decompensation, "strenuous" gymnastics and riding proved to be the most effective. Reducing workplace-base sitting has been shown to address sedentary behaviour in workplace. However, there are few interventions that are cost-effective in reducing occupational sitting time. A study estimated the annual value of nature-based PA conducted in England in 2019 in terms of avoided healthcare and societal costs of six non-communicable diseases (ischaemic heart disease, ischaemic stroke, type 2 diabetes, colon cancer, breast cancer and major depressive disorder) at £108.7million.
Physical activities during leisure and clusters of different forms of meaningfulness
Different forms of physical activity in leisure time can be divided into different clusters of activities that have a common denominator in the form of type of meaningfulness, se model to the right (Lundvall & Schantz 2013).
These separate forms of meaningfulness consist of (i) competition and championship, (ii) nature encounters, (iii) aesthetic-expressive, (iv) fitness gymnastics and play, (v) everyday exercise and (vi) five different basic forms of physical training (aerobic, anaerobic, strength, flexibility and coordination training).
How these different clusters have been treated over time from 1813 to today in a context of teacher training for physical education in the Swedish school system has been described by the Swedish professors in human movement science Suzanne Lundvall & Peter Schantz (2013).
Recommendations for physical activity (including sleep and sedentary behavior)
Global recommendations
The World Health Organization recommend the following:
Adults aged 18–64
1. Adults aged 18–64 should do at least 150 minutes of moderate-intensity aerobic physical activity throughout the week or do at least 75 minutes of vigorous-intensity aerobic physical activity throughout the week or an equivalent combination of moderate- and vigorous-intensity activity.
2. Aerobic activity should be performed in bouts of at least 10 minutes duration.
3. For additional health benefits, adults should increase their moderate-intensity aerobic physical activity to 300 minutes per week, or engage in 150 minutes of vigorous-intensity aerobic physical activity per week, or an equivalent combination of moderate- and vigorous-intensity activity.
4. Muscle-strengthening activities should be done involving major muscle groups on 2 or more days a week.
Adults aged 65+
1. Adults aged 65 years and above should do at least 150 minutes of moderate-intensity aerobic physical activity throughout the week or do at least 75 minutes of vigorous-intensity aerobic physical activity throughout the week or an equivalent combination of moderate- and vigorous-intensity activity.
2. Aerobic activity should be performed in bouts of at least 10 minutes duration.
3. For additional health benefits, adults aged 65 years and above should increase their moderate-intensity aerobic physical activity to 300 minutes per week, or engage in 150 minutes of vigorous-intensity aerobic physical activity per week, or an equivalent combination of moderate-and vigorous-intensity activity.
4. Adults of this age group, with poor mobility, should perform physical activity to enhance balance and prevent falls on 3 or more days per week.
5. Muscle-strengthening activities should be done involving major muscle groups, on 2 or more days a week.
6. When adults of this age group cannot do the recommended amounts of physical activity due to health conditions, they should be as physically active as their abilities and conditions allow.
Children and Adolescents aged 5–17
1. Children and youth aged 5–17 should accumulate at least 60 minutes of moderate- to vigorous-intensity physical activity daily.
2. Amounts of physical activity greater than 60 minutes provide additional health benefits.
Country-level recommendations
Australia, New Zealand, the United Kingdom, Canada and the United States are among the countries that have issued physical activity recommendations.
Predictors of physical activity levels
The amount of physical activity conducted by a population—and by extension the proportion of that population reaching guidelines or other specified thresholds—is dictated by several factors including demographics (e.g., age, sex, ethnicity), population health status, cultural aspects, and the state of the environment itself (e.g. infrastructure that affords physical activity). Demographic groups can also intersect, increasing risk to individuals who are both female and socially disadvantaged for example.
Studies have shown that as the availability of natural environments (e.g., parks, woodlands, inland waters, coasts) increases, more leisure-time physical activity such as walking and cycling are reported. Meteorological conditions have been found to predict physical activity differently in different types of environment. For example, in a large population-based study in England, higher air temperatures and lower wind speeds were associated with increased physical activity.
Globally, in 2016, according to a pooled analysis of 298 population-based surveys, around 81% of students aged 11–17 years were insufficiently physically active. The region with the highest prevalence of insufficient activity in 2016 was high-income Asia Pacific.
As a health indicator
Physical activity, qualified in the form of a physical activity vital sign (PAVS) metric, has been proposed as a screening tool in primary care diagnostics. It has been suggested to correspond with BMI and chronic disease, when coupled with demographic information as well as a tool for identifying patients who do not meet certain physical activity guidelines. Generally, this metric is evaluated by a self-reported medical questionnaire, which can significantly affect the validity and applicability of a PAVS in clinical treatment determination.
| Biology and health sciences | Physical fitness | Health |
2741870 | https://en.wikipedia.org/wiki/Level%20%28optical%20instrument%29 | Level (optical instrument) | A level is an optical instrument used to establish or verify points in the same horizontal plane in a process known as levelling. It is used in conjunction with a levelling staff to establish the relative height or levels (the vertical separation) of objects or marks. It is widely used in surveying and construction to measure height differences and to transfer, measure, and set heights of known objects or marks.
It is also known as a surveyor's level, builder's level, dumpy level or the historic "Y" level. It operates on the principle of establishing a visual level relationship between two or more points, for which an inbuilt optical telescope and a highly accurate bubble level are used to achieve the necessary accuracy. Traditionally the instrument was completely adjusted manually to ensure a level line of sight, but modern automatic versions self-compensate for slight errors in the coarse levelling of the instrument, and are thereby quicker to use.
The optical level should not be confused with a theodolite, which can also measure angles in the vertical plane.
Description
The complete unit is normally mounted on a tripod, and the telescope can freely rotate 360° in a horizontal plane. The surveyor adjusts the instrument's level by coarse adjustment of the tripod legs and fine adjustment using three precision levelling screws on the instrument to make the rotational plane horizontal. The surveyor does this with the use of a bull's eye level built into the instrument mount.
The surveyor looks through the eyepiece of the telescope while an assistant holds a vertical level staff which is graduated in inches or centimeters. The level staff is placed with its foot on the point for which the level measurement is required. The telescope is rotated and focused until the level staff is plainly visible in the crosshairs. In the case of a tilting level, the fine level adjustment is made by an altitude screw, using a high accuracy bubble level fixed to the telescope. This can be viewed by a mirror whilst adjusting, or the ends of the bubble in a "split bubble" display can be viewed within the telescope. This also allows assurance of the accurate level of the telescope whilst the sight is being taken. However, in the case of an automatic level, altitude adjustment is done automatically by a suspended prism due to gravity, as long as the coarse levelling of the instrument base is accurate within certain limits.
When level, the staff graduation readings at the crosshairs and stadia marks are recorded, and an identifying mark or marker placed where the level staff rested on the object or position being surveyed.
Invention
In 1832, English civil engineer and inventor William Gravatt, who was commissioned to examine a scheme for the South Eastern Railway's route from London to Dover, became frustrated with the slow and cumbersome operation of the "Y" level during the survey work, and devised the more transportable, easier-to-use "dumpy" level, so called because of its shorter appearance.
The telescope of the historic "Y" level is held in two brass arms, which are part of the mount and the telescope could be easily removed to allow sighting reversal though 180 degrees or an axial rotation of the telescope; both to compensate for optical collimation errors. Because the telescope is not fixed to the level adjusting mechanism, the "Y" instrument is assembled and disassembled for each sighting station. However, the dumpy level is permanently secured to its two support arms and the levelling mechanism, thereby reducing measurement uncertainty and considerably reducing the time taken to set up the instrument. The dumpy uses the same basic principle of level sighting.
Survey operation
After careful setup of the level, the height of the crosshairs is determined by either sighting from a known benchmark with known height determined by a previous survey or an arbitrary point with an assumed height is used.
Sighting is done with an assistant surveyor who holds a graduated staff vertical at the point under measurement. The surveyor rotates the telescope until the graduated staff is in the crosshairs and records the reading. This is repeated for all sightings from that datum. Should the instrument be moved to another position within sighting distance, it is re-levelled, and a sighting taken of a known level in the previous survey. This relates any new levels to the previous levels.
Variants
The Y level or wye level is the oldest and bulkiest of the older style optical instruments. A low-powered telescope is placed in a pair of clamp mounts, and the instrument then leveled using a spirit level, which is mounted parallel to the main telescope.
The term dumpy level (also builder's level) endures despite the evolution in design. They can be manual or automatic, the latter being much quicker to set up.
A tilting level is a variant which has a precision vertical adjustment screw which tilts both the telescope and the high accuracy bubble level attached to it to make them level. This reduces the complete reliance on the levelling accuracy of the instruments' bottom mount, and the "split bubble" display gives additional assurance that the telescope is level whilst taking the sight. This allows faster operation as the bottom mount need not be truly level, though it will introduce a slight error as the vertical axis of the mount is not completely coincident with the telescope centre. The split bubble works by displaying half of both ends of the bubble side by side in the telescope, and when the curved ends are aligned it is level.
An automatic level, self-levelling level, or builder's auto level includes an internal compensator mechanism (a swinging prism) that, when set close to level, automatically removes any remaining variation. This reduces the need to set the instrument base truly level, as with a dumpy level. Self-levelling instruments are the preferred instrument on building sites, construction, and during surveying due to ease of use and rapid setup time.
A digital electronic level is also set level on a tripod and reads a bar-coded staff using electronic laser methods. The height of the staff where the level beam crosses the staff is shown on a digital display. This type of level removes interpolation of graduation by a person, thus removing a source of error and increasing accuracy. During night time, the dumpy level is used in conjunction with an auto cross laser for accurate scale readings.
A transit level also has the ability to measure both the altitude and azimuth of a target object with respect to a reference in the horizontal plane. The instrument is rotated to sight the target, and the vertical and horizontal angles are read off calibrated scales
In popular culture
In the first chapter of Thomas Hardy's 1887 novel The Woodlanders, the narrator states, "He knew every subtle incline of the ten miles of ground between Abbot's Cernel and Sherton—the market town to which he journeyed—as accurately as any surveyor could have learnt it by a Dumpy level."
In the online game World of Warcraft, there is a quest in Wetlands given by Surveyor Thurdan to retrieve his lost dumpy level. He even comments on the name, saying, "I didn't name the bloody thing, alright? Go look it up!"
| Technology | Surveying tools | null |
2743476 | https://en.wikipedia.org/wiki/Atmospheric%20escape | Atmospheric escape | Atmospheric escape is the loss of planetary atmospheric gases to outer space. A number of different mechanisms can be responsible for atmospheric escape; these processes can be divided into thermal escape, non-thermal (or suprathermal) escape, and impact erosion. The relative importance of each loss process depends on the planet's escape velocity, its atmosphere composition, and its distance from its star. Escape occurs when molecular kinetic energy overcomes gravitational energy; in other words, a molecule can escape when it is moving faster than the escape velocity of its planet. Categorizing the rate of atmospheric escape in exoplanets is necessary to determining whether an atmosphere persists, and so the exoplanet's habitability and likelihood of life.
Thermal escape mechanisms
Thermal escape occurs if the molecular velocity due to thermal energy is sufficiently high. Thermal escape happens at all scales, from the molecular level (Jeans escape) to bulk atmospheric outflow (hydrodynamic escape).
Jeans escape
One classical thermal escape mechanism is Jeans escape, named after British astronomer Sir James Jeans, who first described this process of atmospheric loss. In a quantity of gas, the average velocity of any one molecule is measured by the gas's temperature, but the velocities of individual molecules change as they collide with one another, gaining and losing kinetic energy. The variation in kinetic energy among the molecules is described by the Maxwell distribution. The kinetic energy (), mass (), and velocity () of a molecule are related by . Individual molecules in the high tail of the distribution (where a few particles have much higher speeds than the average) may reach escape velocity and leave the atmosphere, provided they can escape before undergoing another collision; this happens predominantly in the exosphere, where the mean free path is comparable in length to the pressure scale height. The number of particles able to escape depends on the molecular concentration at the exobase, which is limited by diffusion through the thermosphere.
Three factors strongly contribute to the relative importance of Jeans escape: mass of the molecule, escape velocity of the planet, and heating of the upper atmosphere by radiation from the parent star. Heavier molecules are less likely to escape because they move slower than lighter molecules at the same temperature. This is why hydrogen escapes from an atmosphere more easily than carbon dioxide. Second, a planet with a larger mass tends to have more gravity, so the escape velocity tends to be greater, and fewer particles will gain the energy required to escape. This is why the gas giant planets still retain significant amounts of hydrogen, which escape more readily from Earth's atmosphere. Finally, the distance a planet orbits from a star also plays a part; a close planet has a hotter atmosphere, with higher velocities and hence, a greater likelihood of escape. A distant body has a cooler atmosphere, with lower velocities, and less chance of escape.
Hydrodynamic escape
An atmosphere with high pressure and temperature can also undergo hydrodynamic escape. In this case, a large amount of thermal energy, usually through extreme ultraviolet radiation, is absorbed by the atmosphere. As molecules are heated, they expand upwards and are further accelerated until they reach escape velocity. In this process, lighter molecules can drag heavier molecules with them through collisions as a larger quantity of gas escapes. Hydrodynamic escape has been observed for exoplanets close to their host star, including the hot Jupiter HD 209458b.
Non-thermal (suprathermal) escape
Escape can also occur due to non-thermal interactions. Most of these processes occur due to photochemistry or charged particle (ion) interactions.
Photochemical escape
In the upper atmosphere, high energy ultraviolet photons can react more readily with molecules. Photodissociation can break a molecule into smaller components and provide enough energy for those components to escape. Photoionization produces ions, which can get trapped in the planet's magnetosphere or undergo dissociative recombination. In the first case, these ions may undergo escape mechanisms described below. In the second case, the ion recombines with an electron, releases energy, and can escape.
Sputtering escape
Excess kinetic energy from the solar wind can impart sufficient energy to eject atmospheric particles, similar to sputtering from a solid surface. This type of interaction is more pronounced in the absence of a planetary magnetosphere, as the electrically charged solar wind is deflected by magnetic fields, which mitigates the loss of atmosphere.
Charge exchange escape
Ions in the solar wind or magnetosphere can charge exchange with molecules in the upper atmosphere. A fast-moving ion can capture the electron from a slow atmospheric neutral, creating a fast neutral and a slow ion. The slow ion is trapped on the magnetic field lines, but the fast neutral can escape.
Polar wind escape
Atmospheric molecules can also escape from the polar regions on a planet with a magnetosphere, due to the polar wind. Near the poles of a magnetosphere, the magnetic field lines are open, allowing a pathway for ions in the atmosphere to exhaust into space. The ambipolar electric field accelerates any ions in the ionosphere, launching along these lines.
Impact erosion
The impact of a large meteoroid can lead to the loss of atmosphere. If a collision is sufficiently energetic, it is possible for ejecta, including atmospheric molecules, to reach escape velocity.
In order to have a significant effect on atmospheric escape, the radius of the impacting body must be larger than the scale height. The projectile can impart momentum, and thereby facilitate escape of the atmosphere, in three main ways: (a) the meteoroid heats and accelerates the gas it encounters as it travels through the atmosphere, (b) solid ejecta from the impact crater heat atmospheric particles through drag as they are ejected, and (c) the impact creates vapor which expands away from the surface. In the first case, the heated gas can escape in a manner similar to hydrodynamic escape, albeit on a more localized scale. Most of the escape from impact erosion occurs due to the third case. The maximum atmosphere that can be ejected is above a plane tangent to the impact site.
Dominant atmospheric escape and loss processes in the Solar System
Earth
Atmospheric escape of hydrogen on Earth is due to charge exchange escape (~60–90%), Jeans escape (~10–40%), and polar wind escape (~10–15%), currently losing about 3 kg/s of hydrogen. The Earth additionally loses approximately 50 g/s of helium primarily through polar wind escape. Escape of other atmospheric constituents is much smaller. A Japanese research team in 2017 found evidence of a small number of oxygen ions on the moon that came from the Earth.
In 1 billion years, the Sun will be 10% brighter than it is now, making it hot enough on Earth to dramatically increase the water vapor in the atmosphere where solar ultraviolet light will
dissociate H2O, allowing it to gradually escape into space until the oceans dry up
Venus
Recent models indicate that hydrogen escape on Venus is almost entirely due to suprathermal mechanisms, primarily photochemical reactions and charge exchange with the solar wind. Oxygen escape is dominated by charge exchange and sputtering escape. Venus Express measured the effect of coronal mass ejections on the rate of atmospheric escape of Venus, and researchers found a factor of 1.9 increase in escape rate during periods of increased coronal mass ejections compared with calmer space weather.
Mars
Primordial Mars also suffered from the cumulative effects of multiple small impact erosion events, and recent observations with MAVEN suggest that 66% of the 36Ar in the Martian atmosphere has been lost over the last 4 billion years due to suprathermal escape, and the amount of CO2 lost over the same time period is around 0.5 bar or more.
The MAVEN mission has also explored the current rate of atmospheric escape of Mars. Jeans escape plays an important role in the continued escape of hydrogen on Mars, contributing to a loss rate that varies between 160 - 1800 g/s. Jeans escape of hydrogen can be significantly modulated by lower atmospheric processes, such as gravity waves, convection, and dust storms. Oxygen loss is dominated by suprathermal methods: photochemical (~1300 g/s), charge exchange (~130 g/s), and sputtering (~80 g/s) escape combine for a total loss rate of ~1500 g/s. Other heavy atoms, such as carbon and nitrogen, are primarily lost due to photochemical reactions and interactions with the solar wind.
Titan and Io
Saturn's moon Titan and Jupiter's moon Io have atmospheres and are subject to atmospheric loss processes. They have no magnetic fields of their own, but orbit planets with powerful magnetic fields, which protects a given moon from the solar wind when its orbit is within the bow shock. However Titan spends roughly half of its orbital period outside of the bow-shock, subjected to unimpeded solar winds. The kinetic energy gained from pick-up and sputtering associated with the solar winds increases thermal escape throughout the orbit of Titan, causing neutral hydrogen to escape. The escaped hydrogen maintains an orbit following in the wake of Titan, creating a neutral hydrogen torus around Saturn. Io, in its orbit around Jupiter, encounters a plasma cloud. Interaction with the plasma cloud induces sputtering, kicking off sodium particles. The interaction produces a stationary banana-shaped charged sodium cloud along a part of the orbit of Io.
Observations of exoplanet atmospheric escape
Studies of exoplanets have measured atmospheric escape as a means of determining atmospheric composition and habitability. The most common method is Lyman-alpha line absorption. Much as exoplanets are discovered using the dimming of a distant star's brightness (transit), looking specifically at wavelengths corresponding to hydrogen absorption describes the amount of hydrogen present in a sphere around the exoplanet. This method indicates that the hot Jupiters HD209458b and HD189733b and Hot Neptune GJ436b are experiencing significant atmospheric escape.
In 2018 it was discovered with the Hubble Space Telescope that atmospheric escape can also be measured with the 1083 nm Helium triplet. This wavelength is much more accessible from ground-based high-resolution spectrographs, when compared to the ultraviolet Lyman-alpha lines. The wavelength around the helium triplet has also the advantage that it is not severely affected by interstellar absorption, which is an issue for Lyman-alpha. Helium has on the other hand the disadvantage that it requires knowledge about the hydrogen-helium ratio to model the mass-loss of the atmosphere. Helium escape was measured around many giant exoplanets, including WASP-107b, WASP-69b and HD 189733b. It has also been detected around some mini-Neptunes, such as TOI-560 b and HD 63433 c.
Other atmospheric loss mechanisms
Sequestration is not a form of escape from the planet, but a loss of molecules from the atmosphere and into the planet. It occurs on Earth when water vapor condenses to form rain or glacial ice, when carbon dioxide is sequestered in sediments or cycled through the oceans, or when rocks are oxidized (for example, by increasing the oxidation states of ferric rocks from Fe2+ to Fe3+). Gases can also be sequestered by adsorption, where fine particles in the regolith capture gas which adheres to the surface particles.
| Physical sciences | Atmosphere: General | Earth science |
25563683 | https://en.wikipedia.org/wiki/Spanish%20garden | Spanish garden | A traditional Spanish garden is a style of garden or designed landscape developed in historic Spain. Especially in the United States, the term tends to be used for a garden design style with a formal arrangement that evokes, usually not very precisely, the sort of plan and planting developed in southern Spain, incorporating principles and elements from precedents in ancient Persian gardens, Roman gardens and Islamic gardens, and the great Moorish gardens (historically known as riyads) of the Al-Andalus era on the Iberian Peninsula.
In other parts of Spain, public parks and large gardens have been more influenced by the Italian garden, French formal garden, and even the English landscape garden. Spain has a variety of climatic conditions, especially in altitude and rainfall, and modern Spanish gardens are very varied accordingly. Spanish urban housing has long had more apartments than small houses, and the small houses have traditionally lacked front garden, with not that much to the rear either, often just a paved patio with small beds by the walls, and space for plants in pots. Until recently, "full" gardens were mostly found in the country or very large urban houses, but some modern suburban developments have gardens closer to those of northern Europe and North America.
Traditions
Traditionally, the paradise garden is interpreted with a central cross axis, in the four cardinal directions, with long ponds or water channels (a rill or stylized qanat) where water reflects and flows, set in a walled courtyard. The remaining quadrants often had fruit trees and fragrant plants. Thus, characteristic sensory experiences are refreshing coolness, humidity, sounds, greenery, and fragrance. This type of garden is compatible with the Spanish climate of sun and heat. Provisions for shade are given with the use of arcades, pergolas, trellising, and garden pavilions. Ceramic elements and tiles are often used: in water features; for structural, decorative, and seating elements; and as paving; with solid fields, embellishments and accents; and in pottery. A clarity from the symmetrical simplicity often results.
Historical design eras
Spain has a long tradition of making gardens. Significant gardens were made by:
immigrants from the Carthaginian and Roman Empires; for example, the Palmeral of Elche in Alicante
nobility, Christians in the Spanish Medieval period
Islamic rulers and artisans of Al-Andalus, the Moorish Iberian Peninsula or Spanish territories, especially in present-day Andalusia in Southern Spain; for example, the Alhambra, Generalife in Granada.
post-Reconquista Mudéjar design and artisans; for example, the Alcázar of Seville.
catholic monarchs during the Spanish Renaissance, Spanish Gothic, and Spanish Baroque periods; for example, the Royal Palace of La Granja de San Ildefonso.
landowning and business dynasties during the Romantic and Modern periods; for example, Park Güell.
civic projects and expositions; for example, Maria Luisa Park and Plaza de España in Seville.
In 21st century Spain, gardens are designed by garden and landscape designers, horticulturalists, artists, architects, and landscape architects; for example, the Olympic village public outdoor spaces for the 1992 Barcelona Olympics or the public spaces for the Universal Exposition of Seville Expo '92.
Many historic gardens are protected by a heritage designation, Jardín histórico.
Gallery
| Technology | Buildings and infrastructure | null |
523217 | https://en.wikipedia.org/wiki/Atacama%20Large%20Millimeter%20Array | Atacama Large Millimeter Array | The Atacama Large Millimeter/submillimeter Array (ALMA) is an astronomical interferometer of 66 radio telescopes in the Atacama Desert of northern Chile, which observe electromagnetic radiation at millimeter and submillimeter wavelengths. The array has been constructed on the elevation Chajnantor plateau – near the Llano de Chajnantor Observatory and the Atacama Pathfinder Experiment. This location was chosen for its high elevation and low humidity, factors which are crucial to reduce noise and decrease signal attenuation due to Earth's atmosphere. ALMA provides insight on star birth during the early Stelliferous era and detailed imaging of local star and planet formation.
ALMA is an international partnership amongst Europe, the United States, Canada, Japan, South Korea, Taiwan, and Chile. Costing about US$1.4 billion, it is the most expensive ground-based telescope in operation. ALMA began scientific observations in the second half of 2011 and the first images were released to the press on 3 October 2011. The array has been fully operational since March 2013.
Overview
The initial ALMA array is composed of 66 high-precision antennae, and operates at wavelengths of 3.6 to 0.32 millimeters (31 to 1000 GHz). The array has much higher sensitivity and higher resolution than earlier submillimeter telescopes such as the single-dish James Clerk Maxwell Telescope or existing interferometer networks such as the Submillimeter Array or the Institut de Radio Astronomie Millimétrique (IRAM) Plateau de Bure facility.
The antennae can be moved across the desert plateau over distances from 150 m to 16 km, which will give ALMA a powerful variable "zoom", similar in its concept to that employed at the centimeter-wavelength Very Large Array (VLA) site in New Mexico, United States.
The high sensitivity is mainly achieved through the large numbers of antenna dishes that make up the array.
The telescopes were provided by the European, North American and East Asian partners of ALMA. The American and European partners each provided twenty-five 12-meter diameter antennae, for a subtotal of fifty antennae, that compose the main array. The participating East Asian countries are contributing 16 antennae (four 12-meter diameter and twelve 7-meter diameter antennae) in the form of the Atacama Compact Array (ACA), which is part of the enhanced ALMA.
By using smaller antennae than the main ALMA array, larger fields of view can be imaged at a given frequency using ACA. Placing the antennae closer together enables the imaging of sources of larger angular extent. The ACA works together with the main array in order to enhance the latter's wide-field imaging capability.
History
ALMA has its conceptual roots in three astronomical projects: the Millimeter Array (MMA) of the United States, the Large Southern Array (LSA) of Europe, and the Large Millimeter Array (LMA) of Japan.
The first step toward the creation of what would become ALMA came in 1997, when the National Radio Astronomy Observatory (NRAO) and the European Southern Observatory (ESO) agreed to pursue a common project that merged the MMA and LSA. The merged array combined the sensitivity of the LSA with the frequency coverage and superior site of the MMA. ESO and NRAO worked together in technical, science, and management groups to define and organise a joint project between the two observatories with participation by Canada and Spain (the latter became a member of ESO later).
A series of resolutions and agreements led to the choice of "Atacama Large Millimeter Array", or ALMA, as the name of the new array in March 1999 and the signing of the ALMA Agreement on 25 February 2003, between the North American and European parties. ("Alma" means "soul" in Spanish and "learned" or "knowledgeable" in Arabic.) Following mutual discussions over several years, the ALMA Project received a proposal from the National Astronomical Observatory of Japan (NAOJ) whereby Japan would provide the ACA (Atacama Compact Array) and three additional receiver bands for the large array, to form Enhanced ALMA. Further discussions between ALMA and NAOJ led to the signing of a high-level agreement on 14 September 2004 that makes Japan an official participant in Enhanced ALMA, to be known as the Atacama Large Millimeter/submillimeter Array. A groundbreaking ceremony was held on November 6, 2003 and the ALMA logo was unveiled.
During an early stage of the planning of ALMA, it was decided to employ ALMA antennae designed and constructed by known companies in North America, Europe, and Japan, rather than using one single design. This was mainly for political reasons. Although very different approaches have been chosen by the providers, each of the antenna designs appears to be able to meet ALMA's stringent requirements. The components designed and manufactured across Europe were transported by specialist aerospace and astrospace logistics company Route To Space Alliance, 26 in total which were delivered to Antwerp for onward shipment to Chile.
Funding
ALMA was initially a 50-50 collaboration between the National Radio Astronomy Observatory and European Southern Observatory (ESO) and later extended with the help of the other Japanese, Taiwanese, and Chilean partners. ALMA is the largest and most expensive ground-based astronomical project, costing between US$1.4 and 1.5 billion. (However, various space astronomy projects including the Hubble Space Telescope, the James Webb Space Telescope, and several major planet probes have cost considerably more).
Partners
European Southern Observatory and the European Regional Support Centre
National Science Foundation via the National Radio Astronomy Observatory and the North American ALMA Science Center
National Research Council of Canada
National Astronomical Observatory of Japan (NAOJ) under the National Institutes of Natural Sciences (NINS)
ALMA-Taiwan at the Academia Sinica Institute of Astronomy & Astrophysics (ASIAA)
Republic of Chile
Construction
The complex was built primarily by European, U.S., Japanese, and Canadian companies and universities. Three prototype antennae have undergone evaluation at the Very Large Array since 2002.
General Dynamics C4 Systems and its SATCOM Technologies division was contracted by Associated Universities, Inc. to provide twenty-five of the 12 m antennae, while European manufacturer Thales Alenia Space provided the other twenty-five principal antennae (in the largest-ever European industrial contract in ground-based astronomy). Japan's Mitsubishi Electric was contracted to assemble NAOJ's 16 antennae. The antennae were delivered to the site from December 2008 to September 2013.
Transporting the antennae
Transporting the 115 tonne antennae from the Operations Support Facility at 2900 m altitude to the site at 5000 m,
or moving antennae around the site to change the array size, presents enormous challenges; as portrayed in the television documentary Monster Moves: Mountain Mission. The solution chosen is to use two custom 28-wheel self-loading heavy haulers. The vehicles were made by in Germany and are 10 m wide, 20 m long and 6 m high, weighing 130 tonnes. They are powered by twin turbocharged 500 kW Diesel engines.
The transporters, which feature a driver's seat designed to accommodate an oxygen tank to aid breathing the thin high-altitude air, place the antennae precisely on the pads. The first vehicle was completed and tested in July 2007. Both transporters were delivered to the ALMA Operations Support Facility (OSF) in Chile on 15 February 2008.
On 7 July 2008, an ALMA transporter moved an antenna for the first time, from inside the antenna assembly building (Site Erection Facility) to a pad outside the building for testing (holographic surface measurements).
During Autumn 2009, the first three antennae were transported one-by-one to the Array Operations Site. At the end of 2009, a team of ALMA astronomers and engineers successfully linked three antennae at the elevation observing site thus finishing the first stage of assembly and integration of the fledgling array. Linking three antennae allows corrections of errors that can arise when only two antennae are used, thus paving the way for precise, high-resolution imaging. With this key step, commissioning of the instrument began 22 January 2010.
On 28 July 2011, the first European antenna for ALMA arrived at the Chajnantor plateau, 5,000 meters above sea level, to join 15 antennae already in place from the other international partners. This was the number of antennae specified for ALMA to begin its first science observations, and was therefore an important milestone for the project. In October 2012, 43 of the 66 antennae had been set up.
Scientific results
Images from initial testing
By the summer of 2011, sufficient telescopes were operational during the extensive program of testing prior to the Early Science phase for the first images to be captured. These early images gave a first glimpse of the potential of the new array that will produce much better quality images in the future as the scale of the array continues to increase.
The target of the observation was a pair of colliding galaxies with dramatically distorted shapes, known as the Antennae Galaxies. Although ALMA did not observe the entire galaxy merger, the result is the best submillimeter-wavelength image ever made of the Antennae Galaxies, showing the clouds of dense cold gas from which new stars form, which cannot be seen using visible light.
Comet studies
On 11 August 2014, astronomers released studies, using the Atacama Large Millimeter/submillimeter Array (ALMA) for the first time, that detailed the distribution of HCN, HNC, H2CO, and dust inside the comae of comets C/2012 F6 (Lemmon) and C/2012 S1 (ISON).
Planetary formation
An image of the protoplanetary disc surrounding HL Tauri (a very young T Tauri star in the constellation Taurus) was made public in 2014, showing a series of concentric bright rings separated by gaps, indicating protoplanet formation. , most theories did not expect planetary formation in such a young (100,000-1,000,000-year-old) system, so the new data spurred renewed theories of protoplanetary development. One theory suggests that the faster accretion rate might be due to the complex magnetic field of the protoplanetary disc.
Event Horizon Telescope
ALMA participated in the Event Horizon Telescope project, which produced the first direct image of a black hole, published in 2019.
Phosphine in the atmosphere of Venus
ALMA participated in the claimed detection of phosphine, a biomarker, in the air of Venus. As no known non-biological source of phosphine on Venus could produce phosphine in the concentrations detected, this would have indicated the presence of biological organisms in the atmosphere of Venus. Later reanalyses cast doubt on the detection, although later analyses confirmed the results. The detection remains controversial, and is awaiting additional measurements.
Global collaboration
The Atacama Large Millimeter/submillimeter Array (ALMA), an international astronomy facility, is a partnership of Europe, North America and East Asia in cooperation with the Republic of Chile. ALMA is funded in Europe by the European Southern Observatory (ESO), in North America by the U.S. National Science Foundation (NSF) in cooperation with the National Research Council of Canada (NRC) and the National Science Council of Taiwan (NSC) and in East Asia by the National Institutes of Natural Sciences of Japan (NINS) in cooperation with the Academia Sinica (AS) in Taiwan. ALMA construction and operations are led on behalf of Europe by ESO, on behalf of North America by the National Radio Astronomy Observatory (NRAO), which is managed by Associated Universities, Inc (AUI) and on behalf of East Asia by the National Astronomical Observatory of Japan (NAOJ). The Joint ALMA Observatory (JAO) provides the unified leadership and management of the construction, commissioning and operation of ALMA. Its current director since February 2018 is Sean Dougherty.
ALMA regional centre (ARC)
The ALMA regional centre (ARC) has been designed as an interface between user communities of the major contributors of the ALMA project and the JAO. Activates for operating the ARC have also divided into the three main regions involved (Europe, North America and East Asia). The European ARC (led by ESO) has been further subdivided into ARC-nodes located across Europe in Bonn-Bochum-Cologne, Bologna, Ondřejov, Onsala, IRAM (Grenoble), Leiden and JBCA (Manchester).
The core purpose of the ARC is to assist the user community with the preparation of observing proposals, ensure observing programs meet their scientific goals efficiently, run a help-desk for submitting proposals and observing programs, delivering the data to principal investigators, maintenance of the ALMA data archive, assistance with the calibration of data and providing user feedback.
Project detail
At least 50 antennae of 12 m diameter located at an elevation of 5,000 m at Llano de Chajnantor Observatory, enhanced by a compact array of 4 x 12 m and 12 x 7 m antennae (in 2006, consortium considered whether to build 50 or 64 of the 12 m ones. After a Tough Year, ALMA's Star Begins to Rise at Last High and dry)
Imaging instrument in all atmospheric windows between 350 μm and 10 mm
Array configurations from approximately 150 m to 14 km
Spatial resolution of 10 milliarcseconds (10−7 radians), 10 times better than the Very Large Array (VLA) and 5 times better than the Hubble Space Telescope, but still considerably lower than the resolution achieved with optical and infrared interferometers.
The ability to image sources arcminutes to degrees across at one arcsecond resolution
Velocity resolution under 50 m/s
Faster and more flexible imaging instrument than the Very Large Array
Largest and most sensitive instrument in the world at millimeter and submillimeter wavelengths
Point source detection sensitivity 20 times better than the Very Large Array
Data reduction system will be CASA (Common Astronomy Software Applications) which is a new software package based on AIPS++
Atacama Compact Array
The Atacama Compact Array, ACA, is a subset of 16 closely separated antennae that will greatly improve ALMA's ability to study celestial objects with a large angular size, such as molecular clouds and nearby galaxies. The antennae forming the Atacama Compact Array, four 12-meter antennae and twelve 7-meter antennae, were produced and delivered by Japan. In 2013, the Atacama Compact Array was named the Morita Array after Professor Koh-ichiro Morita, a member of the Japanese ALMA team and designer of the ACA, who died on 7 May 2012 in Santiago.
Work stoppage
In August 2013, workers at the telescope went on strike to demand better pay and working conditions. This is one of the first strikes to affect an astronomical observatory. The work stoppage began after the observatory failed to reach an agreement with the workers' union. After 17 days an agreement was reached providing for reduced schedules and higher pay for work done at high altitude.
In March 2020, ALMA was shut down due to the COVID-19 pandemic. It also delayed the cycle 8 proposal submission deadline and suspended public visits to the site.
On October 29, 2022, ALMA suspended observations due to a cyber attack. Observations were restarted 48 days later, on December 16, 2022.
Project timeline
Gallery
| Technology | Ground-based observatories | null |
523380 | https://en.wikipedia.org/wiki/Tomatillo | Tomatillo | The tomatillo (Physalis philadelphica and Physalis ixocarpa), also known as the Mexican husk tomato, is a plant of the nightshade family bearing small, spherical, and green or green-purple fruit of the same name. Tomatillos originated in Mexico and were cultivated in the pre-Columbian era. A staple of Mexican cuisine, they are eaten raw and cooked in a variety of dishes, particularly salsa verde. The tomatillo is a perennial plant, but is generally grown for agriculture each year as if it were an annual.
Names
The tomatillo (from Nahuatl, ) is also known as husk tomato, Mexican groundcherry, large-flowered tomatillo, or Mexican husk tomato. Some of these names, however, can also refer to other species in the genus Physalis. Other names are Mexican green tomato and miltomate.
In Spanish, it is called tomate de cáscara (husk tomato), tomate de fresadilla (little strawberry tomato), tomate milpero (field tomato), tomate verde (green tomato), tomatillo (Mexico; this term means "little tomato" elsewhere), miltomate (Mexico, Guatemala), farolito (little lantern), or simply tomate (in which case the tomato is called jitomate from Nahuatl ).
The tomatillo genus name Physalis is from New Latin physalis, coined by Linnaeus from Ancient Greek φυσαλλίς (physallís, "bladder, wind instrument"), itself from φυσιόω (physióō, "to puff up, blow up"), (physṓ).
Ixocarpa means "slimy fruit", referencing a sticky or slimy coat often on a tomatillo before it ruptures from the calyx.
Distribution
Tomatillos are native to Central America and Mexico, having a wild growth range from Mexico to Costa Rica. The plant is grown mostly in the Mexican states of Hidalgo and Morelos, and in the highlands of Guatemala where it is known as miltomate. In the United States, tomatillos have been cultivated since 1863, with one dubbed "jamberry" in 1945 and others with the names "Mayan husk tomato" and "jumbo husk tomato". Further distribution occurred in the Bahamas, Puerto Rico, Jamaica, and Florida. By the middle of the 20th century, the plant was further exported to India, Australia, South Africa, and Kenya.
The wild tomatillo and related plants are found everywhere in the Americas except in the far north, with the highest diversity in Mexico. In 2017, scientists reported on their discovery and analysis of Physalis infinemundi, a fossil Physalis found in the Patagonian region of Argentina, dated to 52 million years BP. The finding has pushed back the earliest appearance of the Solanaceae plant family and the Physalis genus of which the tomatillo is a part.
Cultivation
History
Tomatillos were domesticated in Mexico before the coming of Europeans and played an important part in the culture of the Maya and the Aztecs, more important than the tomato. The specific name philadelphica dates to the 18th century.
Production
There is limited information about tomatillo production, even though tomatillos are distributed and grown worldwide as a home-grown garden plant. Tomatillos are mainly cultivated in outdoor fields in Mexico and Guatemala on a large scale. Smaller crops are planted in many parts of the United States. In Mexico, tomatillos are planted within a wide range of altitudes.
Soil and climate requirements
In general, tomatillo plants are tolerant to many different soil conditions. However, they do best in well-drained, sandy, fertile soil conditions with a pH between 5.5 and 7.3. Tomatillo plants are cold sensitive. They grow best at . Below , growth is very poor. Tomatillo plants prefer full sun exposure and warm locations.
Seedbed requirement and sowing
Transplanting is the most common practice for planting tomatillo plants. Transplants are produced in greenhouses or transplant beds. Germination occurs at . Transplanting occurs 6 to 8 weeks after seeding and when the risk of frost is past. Transplants produced indoors need to harden off in a warm, sunny place for a few days before being planted outside. Direct outdoor seeding can only be done if no frost risk exists and soil temperature is higher than . Direct outdoor seeding leads to the shortening of the vegetation period. Due to its branching growing pattern, a single plant requires sufficient growing space. Tomatillos are typically grown in rows apart. Although tomatillo is a perennial plant, overwintering is difficult, so it is normally cultivated as an annual plant.
Fertilization and field management
Tomatillo plants can reach heights of . Due to their rapid and branching growth, it is recommended to stake them. Staking also facilitates later harvesting and prevents the fruit from touching the ground, which reduces damage to fruit and husk. Staking can also reduce disease and slug damage. Fertilization is recommended at a moderate level. An application of of phosphorus is common. Other nutrients and fertilizers (N/ K) may be required depending on soil type and irrigation. For non-commercial production, regular fertilization is recommended. Although tomatillo plants become more drought-tolerant as they age, regular watering is required. Tomatillo plants require of water per week. Water can come from rainfall or irrigation. Irrigation can be managed by drip, sprinkler, furrow, or watering can. Irrigation frequency depends on weather and the crop's growth stage, ranging from once or twice a week to daily during hot weather. Weeds are a serious challenge in tomatillo production and are especially important during the first few weeks. Plastic and organic mulches help to control weeds effectively. Applications of plastic mulches also help to restrict soil water evaporation and modify microclimate, thereby affecting tomatillo growth and yield.
Harvest and postharvest treatment
Tomatillos are harvested when the fruits fill the calyx. This state is normally achieved 65 to 100 days after transplanting. Fruit production continues for 1 to 2 months or until the first frost. Harvesting occurs regularly, typically every day, and is done by hand. A plant produces 60 to 200 fruits within a single growing season, with an average yield of about . Tomatillos can be stored for up to three weeks in a cold and humid environment.
Culinary uses
Tomatillos can be harvested at different stages of ripeness. For salsa verde, harvesting may be done early when the fruit is sour with a light flavor. Tomatillos can be picked later when the fruits are seedier for a sweeter taste. Tomatillos have diverse uses in stews, soups, salads, curries, stirfries, baking, cooking with meats, marmalade, and desserts.
Tomatillos are a key ingredient in fresh and cooked Mexican and Central-American green sauces. The green color and tart flavor are the main culinary contributions of the fruit. Purple and red-ripening cultivars often have a slight sweetness, unlike the green- and yellow-ripening cultivars, so they generally are used in jams and preserves. Like their close relative, the Cape gooseberry, tomatillos have a high pectin content. Another characteristic is that they tend to have a varying degree of a sappy, sticky coating, mostly when used on the green side out of the husk.
Ripe tomatillos keep refrigerated for about two weeks. They keep longer with the husks removed and the fruit refrigerated in sealed plastic bags. They may also be frozen whole or sliced.
Tomatillos can also be dried to enhance the sweetness of the fruit in a way similar to dried cranberries, with a hint of tomato flavor. The tomatillo flavor is used in fusion cuisines for blending flavors from Latin American dishes with those of Europe and North America.
Botany
Description
P. ixocarpa is often confused with P. philadelphica due to morphological similarities and the fact that neither species have had a clear type designation. Physalis ixocarpa and Physalis philadelphica have blue anthers that twist after opening, a yellow corolla with five blue-tinged spots or smudges, and a 10-ribbed calyx filled or burst by the berry. The two species differ in flower size and stigma type.
P. philadelphica grow up to and have few hairs on the stem. The leaves have acute and irregularly separated dents on the side. They are typically about in height, and can be either compact and upright or prostrate with a wider, less dense canopy. The leaves are typically serrated and can be either smooth or pubescent.
Classification
The tomatillo is a member of the genus Physalis, erected by Carl Linnaeus in 1753. Jean-Baptiste de Lamarck described the tomatillo under the name Physlis philadelphica in 1786. Other species, such as P. aeuata and P. violacea were described later. The tomatillo is also often classified as P. ixocarpa Brot. However, P. philadelphica is the most important species economically. The nomenclature for Physalis has changed since the 1950s. P. philadelphica was at one time classified as a variety of P. ixocarpa. Later, the classification of P. ixocarpa was revised under the species of P. philadelphica. Today, the name P. ixocarpa is commonly used for the domestic plant and P. philadelphica for the wild one.
Flower
Flowers come in several colors: white, light green, bright yellow, and sometimes purple. Flowers may or may not have purple spots toward the center of the corolla. The anthers are typically dark purple to pale blue. Tomatillo plants are highly self-incompatible, and two or more plants are needed for proper pollination. Thus, isolated tomatillo plants rarely set fruit.
Fruit
The tomatillo fruit is surrounded by an inedible, paper-like husk formed from the calyx. As the fruit matures, it fills the husk and can split it open by harvest time. The husk turns brown, and the fruit can be ripe in several colors, including yellow, green, or even purple. The freshness and greenness of the husk are quality criteria.Flower types:
Varieties
There are several varieties of tomatillos, with differences in tastes, traits, and ripening colors. Some cultivars include Amarylla, Chupon, Gigante, Green Husk, Mexican, Pineapple, Purple de Milpa, Rio Grande Verde, and Yellow.
Genetic
Self-incompatibility trait
Although self-compatibility is common among wild populations, tomatillos carry self-incompatible traits. The plant, i.e., the fertile hermaphrodite, is not able to produce zygotes after self-pollination occurs. This limits the ability to improve tomatillo production regarding the seed quality and the production of varieties.
The self-compatibility gene is situated in the chromosomes of the tomatillo and is not inherited through cytoplasm. Only heterozygous plants can be self-compatible as the trait is controlled by a dominant gene. Tomatillo can thus produce seeds through self-pollination due to the involvement of self-compatibility traits, but the germination viability is different throughout the produced seeds. This suggests that not only incompatible pollen is involved but also inviability at the seedling stage. A study in 2022 using a commercial cultivar found that it was self-compatible and demonstrated incompatibility only in some of the inter-specific hybrid pollinations that were attempted.
Diseases
Tomatillo is generally a resistant crop as long as its climatic requirements are met. However, as with all crops, mass production brings exposure to pests and diseases. As of 2017, two diseases affecting tomatillos have been documented, namely tomato yellow leaf curl virus and turnip mosaic virus. Symptoms of tomato yellow leaf curl virus, including chlorotic margins and interveinal yellowing, were found in several tomato and tomatillo crops in Mexico and Guatemala in 2006. After laboratory tests, the virus was confirmed. Symptomatic plants were associated with the presence of whiteflies, which were likely the cause of this outbreak.
Turnip mosaic virus was discovered in several tomatillo crops in California in 2011, rendering 2% of commercially grown tomatillo plants unmarketable, with severe stunting and leaf distortion. The green peach aphid is a common pest in California, and since it readily transmits the turnip mosaic virus, this could be a threat to tomatillo production in California.
| Biology and health sciences | Botanical fruits used as culinary vegetables | Plants |
523382 | https://en.wikipedia.org/wiki/AVE | AVE | Alta Velocidad Española (AVE) is a high-speed rail service operated by Renfe, the Spanish State railway company.
The first AVE service was inaugurated in 1992, with the introduction of the first Spanish high-speed railway connecting the cities of Madrid, Córdoba and Seville.
In addition to Renfe's use of the Administrador de Infraestructuras Ferroviarias-managed rail infrastructure in Spain, Renfe offers two AVE services partially in France, connecting respectively Barcelona-Lyon and Madrid-Marseille.
translates to "Spanish High Speed", but the initials are also a play on the word , meaning "bird". AVE trains operate at speeds of up to .
Services
Renfe offers the following AVE services:
Alicante–León via Albacete, Cuenca, Madrid Chamartín, Valladolid and Palencia.
Alicante–Ourense via Albacete, Cuenca, Madrid Chamartín and Zamora.
Barcelona–Granada via Tarragona, Lleida, Zaragoza, Ciudad Real, Puertollano, Córdoba and Antequera.
Barcelona–Málaga via Tarragona, Lleida, Zaragoza, Ciudad Real, Córdoba, Puente Genil-Herrera, and Antequera.
Barcelona–Seville via Tarragona, Lleida, Zaragoza, Ciudad Real, Puertollano and Córdoba (trains with selective stops are also scheduled).
Burgos–Murcia via Valladolid, Segovia, Madrid-Chamartín, Elche and Orihuela.
Gijón–Castellón via Oviedo, Mieres Del Camín, La Pola, León, Palencia, Valladolid, Segovia, Madrid-Chamartín, Cuenca, Valencia and Sagunto.
Gijón–Vinaros, via Oviedo, Mieres Del Camín, La Pola, León, Palencia, Valladolid, Segovia, Madrid-Chamartín, Cuenca, Valencia, Sagunto, Castellón, Benicàssim, Oropesa del Mar and Benicarló (only in summertime).
Huesca–Seville via Tardienta, Zaragoza, Calatayud, Guadalajara, Madrid-Puerta de Atocha and Córdoba
Madrid–A Coruña via Zamora, Ourense and Santiago De Compostela.
Madrid–Alicante via Cuenca, Albacete, and Villena (non stop trains and trains with selective stops are also scheduled).
Madrid–Barcelona via Guadalajara, Calatayud, Zaragoza, Lleida, and Tarragona (non stop trains and trains with selective stops are also scheduled).
Madrid–Castellón via Cuenca, Requena-Utiel and Valencia.
Madrid–Figueres via Guadalajara, Calatayud, Zaragoza, Lleida, Tarragona, Barcelona and Girona (trains are scheduled with selective stops).
Madrid–Gijón via Valladolid, Palencia, León and Oviedo.
Madrid–Granada via Ciudad Real, Puertollano, Córdoba, Puente Genil-Herrera, Antequera and Loja (trains with selective stops are also scheduled).
Madrid–Huesca via Guadalajara, Calatayud, Zaragoza, and Tardienta.
Madrid–León via Segovia, Valladolid and Palencia.
Madrid–Málaga via Ciudad Real, Puertollano, Córdoba, Puente Genil-Herrera, and Antequera (non stop trains and trains with selective stops are also scheduled).
Madrid–Murcia via Elche and Orihuela (some trains are arriving to Alicante and then reversing towards Murcia).
Madrid–Ourense via Zamora.
Madrid–Seville via Ciudad Real, Puertollano, and Córdoba (non stop trains and trains with selective stops are also scheduled).
Madrid–Valencia via Cuenca and Requena-Utiel (non stop trains are also scheduled).
Madrid–Vigo via Zamora, Sanabria, A Gudiña, Ourense, Santiago de Compostela, Vilagarcía de Arousa and Pontevedra (trains with selective stops are also scheduled).
Valencia–Burgos via Requena-Utiel, Cuenca, Madrid Chamartín and Valladolid (trains with selective stops are also scheduled).
Valencia–León via via Requena-Utiel, Cuenca, Madrid-Chamartín, Segovia, Valladolid and Palencia (trains with selective stops are also scheduled).
Valencia–Seville via Cuenca, Ciudad Real, Puertollano, and Córdoba.
International:
Barcelona–Lyon via Girona, Figueres, Perpignan, Narbonne, Montpellier, Nîmes, and Valence.
Madrid–Marseille via Guadalajara, Zaragoza, Tarragona, Barcelona, Girona, Figueres, Perpignan, Narbonne, Béziers, Montpellier, Nîmes, Avignon and Aix-en-Provence.
The central hub of the AVE system is Madrid's Puerta de Atocha, except for the Madrid–Asturias, Madrid–Burgos, Madrid–Galicia, Madrid–Alicante and Madrid-Murcia lines, that terminate at Chamartín station.
Trains
Currently, there are several series of high-speed trains that run the AVE service:
S-100, manufactured by Alstom
S-102, manufactured by Talgo and Bombardier
S-103, manufactured by Siemens, marketed globally under the brand Siemens Velaro
S-106, manufactured by Talgo, marketed globally as Talgo AVRIL.
S-112, manufactured by Talgo and Bombardier
Passenger usage
The still-growing network transported a record 32.4 million passengers in 2023. Though the network length is extensive, it lags in ridership behind comparable high-speed rail systems in Japan, France, Germany, China, Taiwan, and Korea.
Rail infrastructure in Spain and Europe
Rail transport in Spain
High-speed rail in Spain
High-speed rail in Europe
Train categories in Europe
| Technology | High-speed rail | null |
523404 | https://en.wikipedia.org/wiki/Uakari | Uakari | Uakari (, ) is the common name for the New World monkeys of the genus Cacajao. Both the English and scientific names are believed to have originated from indigenous languages.
The uakaris are unusual among New World monkeys in that the tail length (15–18 cm) is substantially less than their head and body length (40–45 cm). Their bodies are covered with long, loose hair but their heads are bald. They have almost no subcutaneous fat, so their bald faces appear almost skull-like. Like their closest relatives the saki monkeys, they have projecting lower incisors. These monkeys have the most striking red facial skin of any primate. Females choose their mates based on how red the male's face is. Evidence suggests that the red facial coloration reflects the health of the primate.
The four species of uakari currently recognized are all found in the north-western Amazon basin. The bald uakari, remarkable for its brilliant scarlet complexion, is found north of the Amazon River, and south of the Japurá River in the Mamirauá Sustainable Development Reserve. The black-headed uakari is found north of the Amazon and south of the Rio Negro. The Neblina uakari is found north of the Rio Negro, west of the Rio Marauiá and east of the Casiquiare canal. The Aracá uakari is currently known only from the Rio Curuduri basin.
They have been observed both in small groups and in larger troops of up to 100. When traveling through the forest they move in the lower branches of the trees, though when foraging they also go up to the canopy. They mostly eat fruit, and unlike other Neotropical frugivores will consume a large amount of unripe fruit for which they have specialised dentition. They also eat flowers, seeds, invertebrates, buds and leaves.
Uakari are found in neotropical Amazon flooded or riparian forests, including Brazil, Colombia, Peru and Venezuela.
A phylogeographic reconstruction found that the concestor of living uakari dates to 1.7 million years ago, in the Solimões River, whence they spread and diversified following intermittent river rearrangements.
Species
There are four species in this genus:
Genus Cacajao
Bald uakari or red uakari, C. calvus
Cacajao calvus calvus
Cacajao calvus ucayalii
Cacajao calvus rubicundus
Cacajao calvus novaesi
Black-headed uakari species group
Golden-backed or black-headed uakari, Cacajao melanocephalus
Aracá uakari, Cacajao ayresi*
Neblina uakari, Cacajao hosomi*
In 2014 Ferrari et al. proposed an alternative taxonomy for the C. melanocephalus group which recognizes the Aracá uakari as a subspecies of the golden-backed uakari, and also recognized Cacajao ouakary as a separate species, whereas current consensus is that C. ouakary is a junior synonym of C. melanocephalus. This revision is not universally accepted.
| Biology and health sciences | New World monkeys | Animals |
523475 | https://en.wikipedia.org/wiki/Korea%20Train%20Express | Korea Train Express | Korea Train eXpress (), often known as KTX (), is South Korea's high-speed rail system, operated by Korail. Construction began on the high-speed line from Seoul to Busan in 1992. KTX services were launched on April 1, 2004.
From Seoul Station the KTX lines radiate with stops at Seoul Station, Yongsan station towards Busan and Gwangju. A new line from Wonju to Gangneung was completed in December 2017 to serve the 2018 Winter Olympics in Pyeongchang.
The current maximum operating speed for trains in regular service is , though the infrastructure is designed for .
The initial rolling stock was based on Alstom's TGV Réseau, and was partly built in Korea. The domestically developed HSR-350x, which achieved in tests, resulted in a second type of high-speed trains now operated by Korail, the KTX-Sancheon, which entered into commercial service in 2010.
The next generation experimental electric multiple unit prototype, HEMU-430X, achieved in 2013, making South Korea the world's fourth country after Japan, France and China to develop a high-speed train running on conventional rail above . It was further developed into commercialised variants, namely KTX-Eum and KTX-Cheongryong, with respective maximum service speeds of and , which entered into KTX services in 2021 and 2024, respectively.
In June 2024, South Korea and Uzbekistan concluded a KRW 270 billion (approximately US$196 million) deal to apply KTX technology in Uzbekistan by supplying high-speed trains and Korail expertise.
History
Origins of the project
The Seoul-Busan axis is Korea's main traffic corridor. In 1982, it represented 65.8% of South Korea's population, a number that grew to 73.3% by 1995, along with 70% of freight traffic and 66% of passenger traffic. With both the Gyeongbu Expressway and Korail's Gyeongbu Line congested as of the late 1970s, the government saw the pressing need for another form of transportation.
The first proposals for a second Seoul-Busan railway line originated from a study prepared between 1972 and 1974 by experts from France's SNCF and the Japan Railway Technical Service on a request from the IBRD. A more detailed 1978-1981 study by KAIST, focusing on the needs of freight transport, also came to the conclusion that separating long-distance passenger traffic on a high-speed passenger railway would be advisable, and it was adopted in the following Korean Five Year Plan.
During the following years, several feasibility studies were prepared for a high-speed line with a Seoul–Busan travel time of 1 hour 30 minutes, which gave positive results. In 1989, following the go-ahead for the project, the institutions to manage its preparation were established: the Gyeongbu High Speed Electric Railway & New International Airport Committee, and the High Speed Electric Railway Planning Department (later renamed HSR Project Planning Board). In 1990, the planned Seoul–Busan travel time was 1 hour 51 minutes, the project was to be completed by August 1998, and costs were estimated at 5.85 trillion South Korean won (₩) in 1988 prices, 4.6 trillion of which were to be spent on infrastructure, the remainder on rolling stock.
As planning progressed, the Korea High Speed Rail Construction Authority (KHSRCA) was established in March 1992 as a separate body with its own budget responsible for the project. In the 1993 reappraisal of the project, the completion date was pushed back to May 2002, and cost estimates grew to ₩10.74 trillion. 82% of the cost increase was due to a 90% increase in unit costs in the construction sector, mostly labour costs but also material costs, and the remainder due to alignment changes. To finance the project, the option of a build-operate-transfer (BOT) franchise was rejected as too risky. Funding included direct government grants (35%), government (10%) and foreign (18%) loans, domestic bond sales (31%) and private capital (6%).
Creation of the system
Start of high-speed line construction
KHSRCA started construction of the Seoul–Busan Gyeongbu high-speed railway (Gyeongbu HSR) on June 30, 1992, on the long section from Cheonan to Daejeon, which was intended for use as test track.
Construction started before the choice of the main technology supplier, thus alignment design was set out to be compatible with all choices. Of the planned line, would be laid on bridges, and another in tunnels. However, plans were changed repeatedly, in particular those for city sections, following disputes with local governments, while construction work suffered from early quality problems. Planned operating speed was also reduced from to the maximum of high-speed trains on the market. Three competitors bid for the supply of the core system, which included the rolling stock, catenary and signalling: consortia led by GEC-Alsthom, today Alstom, one of the builders of France's TGV trains; Siemens, one of the builders of Germany's ICE trains; and Mitsubishi Heavy Industries, one of the builders of Japan's Shinkansen trains. In 1994, the alliance of GEC-Alsthom and its Korean subsidiary Eukorail were chosen as winner.
The technology was almost identical to that found on the high-speed lines of France's TGV system. Track-related design specifications included a design speed of and standard gauge.
Phase 1: Seoul–Daegu and conventional line upgrades
Following the 1997 Asian Financial Crisis, the government decided to realise the Gyeongbu HSR in two phases. In a first phase, two-thirds of the high-speed line between the southwestern suburbs of Seoul and Daegu would be finished by 2004, with trains travelling along the parallel conventional line along the rest of the Seoul–Busan route. The upgrade and electrification of these sections of the Gyeongbu Line was added to the project, and also the upgrade and electrification of the Honam Line from Daejeon to Mokpo, providing a second route for KTX services. The budget for the first phase was set at ₩12,737.7 billion, that for the entire project at ₩18,435.8 billion in 1998 prices. While the share of government contributions remained unchanged, the share of foreign loans, domestic bond sales and private capital changed to 24%, 29% and 2%.
The infrastructure and rolling stock were created in the framework of a technology transfer agreement, which paired up Korean companies with core system supplier Alstom and its European subcontractors for different subsystems. Alstom's part of the project amounted to US$2.1 billion or €1.5 billion.
Well ahead of the opening of the Gyeongbu HSR for regular service, in December 1999, of the test section, later extended to , was finished to enable trials with trains. After further design changes, the high-speed tracks were finished over a length of , with of interconnections to the conventional Gyeongbu Line, including at a short interruption at Daejeon. The high-speed section itself included of viaducts and of tunnels. Conventional line electrification was finished over the across Daegu and on to Busan, the across Daejeon, and the from Daejeon to Mokpo and Gwangju. After 12 years of construction and with a final cost of ₩12,737.7 billion, the initial KTX system with the first phase of the Gyeongbu HSR went into service on April 1, 2004.
Phase 2: Daegu–Busan, extra stations, urban sections
The Daegu–Busan section of the Gyeongbu HSR became a separate project with the July 1998 project revision, with a budget of ₩5,698.1 billion, with funding from the government and private sources by the same ratios as for phase 1. In August 2006, the project was modified to again include the Daejeon and Daegu urban area passages, as well as additional stations along the phase 1 section. For these additions, the budget as well as the government's share of the funding was increased.
Construction started in June 2002. The line, which follows a long curve to the northeast of the existing Gyeongbu Line, includes 54 viaducts with a total length of and 38 tunnels with a total length of . The two largest structures are the Geomjeung Tunnel, under Mount Geumjeong at the Busan end of the line; and the Wonhyo Tunnel, under Mount Cheonseong south-west of Ulsan, which will be the longest and second longest tunnels in Korea once the line is opened.
A long dispute concerning the environmental impact assessment of the Wonhyo Tunnel, which passes under a wetland area, caused delays for the entire project. The dispute gained nationwide and international attention due to the repeated hunger strikes of a Buddhist nun, led to a suspension of works in 2005, and only ended with a supreme court ruling in June 2006. With the exception of the sections across Daejeon and Daegu, the second phase went into service on November 1, 2010. By that time, ₩4,905.7 billion was spent out of a second phase budget, or ₩17,643.4 billion out of the total.
The two sections across the urban areas of Daejeon and Daegu, altogether , will be finished by 2014. As of October 2010, the total cost of the second phase was estimated at ₩7,945.4 billion, that for the entire project at ₩20,728.2 billion. The last element of the original project that was shelved in 1998, separate underground tracks across the Seoul metropolitan area, was re-launched in June 2008, when an initial plan with a long alignment and two new stations was announced.
Further upgrades of connecting conventional lines
The electrification and the completion of the re-alignment and double-tracking of the Jeolla Line, which branches from the Honam Line at Iksan and continues to Suncheon and Yeosu, began in December 2003, with the aim to introduce KTX services in time for the Expo 2012 in Yeosu. The upgrade will allow to raise top speed from . The section of the perpendicular Gyeongjeon Line from Samnangjin, the junction with the Gyeongbu Line near Busan, to Suncheon is upgraded in a similar way, with track doubling, alignment modifications and electrification for . The until Masan was opened on December 15, 2010. The upgrade is to be complete until Jinju by 2012 and Suncheon by 2014. The top speed of the AREX line, Seoul's airport link, is to be raised from for the KTX.
The Ulsan–Gyeongju–Pohang section of the Donghae Line is foreseen for an upgrade in a completely new alignment that circumvents downtown Gyeongju and connects to the Gyeongbu High Speed Railway at Singyeongju Station, allowing for direct KTX access to the two cities. On April 23, 2009, the project was approved by the government and a ground-breaking ceremony was held. The altogether line is slated to be opened in December 2014.
On September 1, 2010, the South Korean government announced a strategic plan to reduce travel times from Seoul to 95% of the country to under 2 hours by 2020. The main new element of the plan is to aim for top speeds of in upgrades of much of the mainline network with view to the introduction of KTX services. The conventional lines under the scope of the plan include the above, already on-going projects, and their extensions along the rest of the southern and eastern coasts of South Korea, lines along the western coast, lines north of Seoul, and the second, more easterly line between Seoul and Busan with some connecting lines.
Further high-speed lines
Honam HSR
Until 2006, the first plans for a second, separate high-speed line from Seoul to Mokpo were developed into the project of a line branching from the Gyeongbu HSR and constructed in two stages, the Honam High Speed Railway (Honam HSR). The budget for the first stage, from the new Osong Station on the Gyongbu HSR to Gwangju·Songjeong Station, was set at ₩8,569.5 billion. The second stage, the remaining to Mokpo, was to be finished by 2017 with a budget of ₩2,002.2 billion. The Osong-Iksan section of the first phase is also intended for use as high-speed test track for rolling stock development, to be fitted with special catenary and instrumented track. The ground-breaking ceremony was held on December 4, 2009. As of September 2010, progress was 9.6% of the project budget then estimated at ₩10,490.1 billion for the first phase, which was due for completion in 2014, while the estimate for the entire line stood at ₩12,101.7 billion.
Suseo HSR
First plans for the Honam HSR foresaw a terminus in Suseo station, southeast Seoul. The branch to Suseo was re-launched as a separate project, the Suseo High Speed Railway (Suseo HSR), in June 2008. Detailed design of the line is underway since September 2010, with opening planned by the end of 2014. For the longer term, new high-speed lines from Seoul to Sokcho on the eastern coast, and a direct branch from the Gyeongbu HSR south to Jinju and further to the coast are under consideration. In conjunction with the award of the 2018 Winter Olympics to PyeongChang in July 2011, KTX service via the eastern coast line was anticipated; the expected travel time there from Seoul is 50 minutes.
Jeju Island
In January 2009, the Korea Transport Institute also proposed a line from Mokpo to Jeju Island, putting Jeju 2 hours 26 minutes from Seoul. The line would include a bridge from Haenam to Bogil Island and a undersea tunnel from Bogil Island to Jeju Island (with a drilling station on Chuja Island), for an estimated cost of US$10 billion. As the proposal was popular with lawmakers from South Jeolla Province, the government is conducting a feasibility study, but the Jeju governor expressed skepticism. The Seoul-Jeju route is the world's busiest air route with 13.7 million passengers (2023).
However, Jeju Gov. Won Hee-ryong opposed this plan since it would ruin the island's identity and make the Jeju economy more dependent on the mainland.
Technology
The shock absorption design absorbs 80 percent of the shock energy when the train crashes. Automatic ventilation is installed to prevent noise from occurring when trains enter and exit the tunnel. Articulated bogies help increase ride comfort and driving safety.
On June 14, 2024, Uzbekistan and South Korea finalized a US$196 million deal for KTX technology to be applied in Uzbekistan. This was the first time KTX technology was exported. As part of the deal, 42 train units capable of going up to were to be supplied for of rail in Uzbekistan. Operations were scheduled to begin in April 2027.
Rolling stock
KTX-I
The initial KTX-I trainsets, also known as simply KTX or as TGV-K, are based on the TGV Réseau, but with several differences. 46 trains were built - the initial twelve in France by Alstom, the remainder in South Korea by Rotem. The 20-car electric multiple units consist of two traction heads, which are powered end cars without passenger compartments, and eighteen articulated passenger cars, of which the two extreme ones have one motorised bogie each. A KTX-I was built to carry up to 935 passengers at a regular top speed of , later increased to .
KTX-Sancheon
For less frequented relations and for operational flexibility, a 2001 study proposed a train created by scaling down the planned commercial version of the HSR-350x, by shortening the train, removing powered bogies from intermediate cars, and lowering top speed.
Hyundai Rotem received orders for altogether 24 such trains, called KTX-II, in three batches from July 2006 to December 2008.
Design speed is , and revenue service speed is . The power electronics uses newer technology than the HSR-350x, and the front is a new design, too. The trainsets, of which two can be coupled together, consist of two traction heads and eight articulated passenger cars, and seat 363 passengers in two classes, with enhanced comfort relative to the KTX-I. The domestic added value of the trains was increased to 87%, compared to 58% for the KTX-I. Imported parts include the pantographs, semiconductors in the power electronics, front design, couplers and final drives.
The train was developed on the basis of the transferred TGV technology, but more advanced technology was used for the new motors, power electronics and additional brake systems, while the passenger cars were made of aluminum to save weight, and the nose was a new design with reduced aerodynamic drag. Test runs were conducted between 2002 and 2008, in the course of which HSR-350x achieved the South Korean rail speed record of on December 16, 2004.
The KTX-II was officially renamed as KTX-Sancheon () after the Korean name of the indigenous fish cherry salmon before the first units started commercial service on March 2, 2010.
However within weeks of its initial launch, mechanical and design flaws began to appear, in some cases causing trains to stop running and forcing passengers to leave the train and walk back to the station, and in one particular case derailing from the tracks on February 11, 2011. Although the trains were designed to be a domestically built replacement for the French built Alstrom trains, due to over 30 malfunctions since March 2, 2010, Korail asked manufacturer Hyundai-Rotem to recall all 19 of the trains in operation after finding cracks in two anchor bands in May 2011.
Following the recall, the KTX-Sancheon trains were put back in service.
In addition to the 24 initial KTX-Sancheon trains, which form the KTX-Sancheon Class 11, new batches have been ordered and delivered since, to provide service on the new Honam, Suseo and Gyeonggang lines. For the opening of the Honam HSR line, 22 trainsets, named Class 12, were delivered ahead of the 2015 opening. In addition, 10 trainsets have been delivered to provide service on the Suseo line, scheduled to open in December 2016 (Class 13), and 15 trainsets (Class 14) have been ordered for the Gyeonggang Line, which opened in late 2017 ahead of the 2018 Winter Olympics.
KTX-Eum
KTX-Eum is South Korea's first high-speed electric multiple unit train, being a commercialised version of the experimental HEMU-430X previously tested by Korail. It has the maximum service speed of and is aimed for serving semi-high-speed railway lines.
The KTX-Eum entered service on Jungang Line on January 4, 2021, operating between electrified section of Cheongnyangni and Andong. It was also introduced on Gangneung Line since August 1, 2021, replacing KTX-Sancheon which would be redeployed to other KTX lines.
A further order of 14 six-car units was placed in December 2016, both orders are to be delivered in 2020–2021.
KTX-Cheongryong
KTX-Cheongryong is a sister train of KTX-Eum, but a trainset consists of eight cars as opposed to six cars. It has the maximum service speed of and is aimed for supplementing trainsets for current high-speed rail services.
In 2016, a contract was concluded between Korail and Hyundai Rotem to build 2 pre-series sets of KTX-Cheongryong. These trainsets entered service on the Gyeongbu high-speed railway and Honam high-speed railway on May 1, 2024.
Another 17 trainsets ordered by Korail are scheduled to be delivered between April 2027 and March 2028.
List of KTX lines
Current lines
Future lines
Defunct lines
Operation
Following a phase of test operation, regular KTX service started on April 1, 2004, with a maximum speed of achieved along the finished sections of the Gyeongbu HSR. In response to frequent passenger complaints regarding speeds on the video display staying just below the advertised 300 mark, operating top speed was raised to on November 26, 2007.
Services
KTX services are grouped according to their route, and within the groups, the stopping pattern changes from train to train. KTX trains not deviating from the Seoul–Busan corridor are operated as the Gyeongbu KTX service. In 2004, the new service cut the route length from , and the fastest trains, serving four stations only, cut the minimum Seoul–Busan travel time from the Saemaul's 4 hours 10 minutes to 2 hours 40 minutes. With the extension of the Gyeongbu HSR, from November 1, 2010, the minimum Seoul–Busan travel time reduced to 2 hours 18 minutes, over a travel distance of . From December 1, 2010, Korail added a pair of non-stop trains with a travel time of 2 hours 8 minutes. Once the sections across Daejeon and Daegu are completed, cutting the Seoul–Busan travel distance to , plans foresee a further improvement of the four-stop travel time to 2 hours and 10 minutes.
Because both KTX and conventional trains in South Korea share a rail gauge (unlike in Japan), KTX trains can run on both networks dramatically increasing the number of destinations served.
Some Gyeongbu KTX services use parts of the conventional line paralleling the high-speed line. From June 2007 until October 2010, some trains left the Gyeongbu HSR between Daejeon and Dongdaegu to serve Gimcheon and Gumi before the opening of an extra station for the two cities on the high-speed line. From November 1, 2010, when most Gyeongbu KTX services began to use the new Daegu–Busan high-speed section, some trains remained on the Gyeongbu Line on that section, and additional trains began to use the Gyeongbu Line on the Seoul–Daejeon section to serve Suwon.
KTX trains using the Gyeongbu HSR only from Seoul to Daejeon and continuing all along the Honam Line are operated as the Honam KTX service. In 2004, the new service with a route length of between Yongsan in Seoul and Mokpo cut minimum travel time from 4 hours 42 minutes to 2 hours 58 minutes. By 2017, this time is to be cut further to 1 hours 46 minutes.
On December 15, 2010, the new Gyeongjeon KTX service started with a minimum travel time of 2 hours 54 minutes over the long route between Seoul and Masan. The service is to be extended to Jinju by 2012. A fourth line, the Jeolla KTX service will connect Seoul to Yeosu in 3 hours 7 minutes from September 2011. From 2014, with the completion of the first phase of the Honam HSR, the travel time is reduced further to 2 hours 25 minutes. From 2015, KTX trains are to reach Pohang from Seoul in 1 hour 50 minutes.
Tickets and seats
Type of seats
KTX offers two classes: First Class and Standard Class. Tickets also specify whether a seat is forward-facing or backward-facing according to the direction of travel. First Class seats are arranged 2+1 across the train and Standard Class seats are configured 2+2. There are special reserved Family seats, which are grouped in four, including 2 forward-facing and 2 backward-facing seats. There are reserved seats and unassigned seats. KTX trains have no restaurant cars or bars, only seat service. From 2006, one car of selected KTX services functions as a moving cinema.
Ticket prices
KTX fares were designed to be about halfway between those for conventional trains and airline tickets. The fare system implemented at the start of service in April 2004 deviated from prices proportional with distance, to favour long-distance trips. On April 25, 2005, fares were selectively reduced for relations under-performing most.
From November 1, 2006, due to rising energy prices, Korail applied an 8-10% fare hike for various train services, including 9.5% for KTX. The price of a Seoul-Busan Standard Class ticket increased to 48,100 won. From July 1, 2007, KTX fares were hiked another 6.5%, while those for the slower Saemaeul and Mugunghwa services on the parallel conventional route were raised by 3.5 percent and 2.5 percent, respectively. However, new reduced weekday and unassigned seat fares were also introduced.
After the November 1, 2010, start of service on the Daegu–Busan section of the Gyeongbu HSR, the fare for KTX trains using the new section was set about 8% higher than for the old route via Miryang, while that for the new services via Suwon was set lower.
Discounts
Korail's standard discounts for children, disabled, seniors and groups apply on KTX trains, too. For frequent travellers, Korail's standard discount cards, which are categorised according to age group, apply with the double of the standard discount rates; while discount cards for business and government agency workers apply with the normal rate; both types of discounts are up to 30%. Season period tickets with discounts of up to 60% can also apply to KTX trains.
Discounts for family seats (37.5%) and backward facing seats (5%) are specific to the KTX. In addition to Korail's small general discounts for tickets purchased in a vending machine, via cell phone or the internet, discounts of 5–20% apply to a limited number of seats on KTX trains when purchased in advance. For travellers who transfer to other long-distance trains towards destinations beyond KTX stops, transfer tickets with 30% discount apply. Korail pays a refund for late KTX trains, which reaches 100% for trains with a delay above one hour.
Korea Rail Pass, a period ticket Korail offers to foreigners, also applies to KTX. For passengers using the Korea-Japan Joint Rail Pass, a joint offer of Korail, Japanese railways and ferry services, the discount on KTX trains is 30%.
Passenger numbers and usage
Forecasts
When the project was launched, KTX was expected to become one of the world's busiest high-speed lines. The first study in 1991 forecast around 200,000 passengers a day in the first year of operation, growing to 330,000 passengers a day twelve years later. In forecasts prepared after the decision to split the project into two phases, the expected first year ridership of Gyeongbu KTX services was reduced by about 40%. With the estimate for the Honam KTX services added to the plan, opening year forecasts ranged between 150,000 and 175,000 passengers a day. Actual initial ridership after the opening of the first phase in 2004 was well short of initial expectations at around half of the final forecast.
In October 2010, before the opening of the second phase, Korail expected ridership to rise from the then current 106,000 to 135,000 passengers a day.
Ridership evolution
KTX was introduced on 1 April 2004. In the first 100 days, daily passenger numbers averaged 70,250, generating an operational revenue of about 2.11 billion won per day, 54% of what was expected. On January 14, 2005, Prime Minister Lee Hae Chan stated that "the launch of KTX was a classic policy failure" due to construction costs significantly above and passenger numbers well below forecasts. However, ridership increased by over a third on the Gyeongbu KTX and over a half on the Honam KTX in two years. Financial break-even was forecast at a ridership level of around 100,000 passengers a day, which was expected by the end of 2006.
The 100 millionth rider was carried after 1116 days of operation on April 22, 2007, when cumulative income stood at 2.78 trillion won. KTX finances moved into the black in 2007. The next year, with revenues equal to US$898 million and costs equal to US$654 million, KTX was Korail's most profitable branch.
By the sixth anniversary in April 2010, KTX trains travelled a total 122.15 million kilometres, carrying 211.01 million passengers. Punctuality gradually improved from 86.7% of trains arriving within 5 minutes of schedule in 2004 to 98.3% in 2009. In 2009, the average daily ridership was 102,700. As of April 2010, the single-day ridership record stood at 178,584 passengers, achieved on January 26, 2009, the Korean New Year.
By the tenth anniversary KTX had travelled a total 240 million kilometres, carrying 414 million passengers.
Market share and effect
The introduction of high-speed services had the strongest effect on long-distance relations with a significant portion of the journey on the high-speed line, like Seoul–Busan: KTX took both the majority of the market and the bulk of rail passengers in the first year already, increasing the total share of rail from around two-fifths to a market dominating two-thirds by 2008. On long-distance relations with significant distances along conventional lines and resulting more modest travel time gains, that is along the Honam Line, the KTX and overall rail market share gain decreases with distance. On medium-distance relations like Seoul–Daejeon, KTX gained market share mostly at the expense of normal rail express services and air traffic, and helped to increase the total share of rail. On short-distance intercity relations line Seoul–Cheonan, due to the modest gains in time and the location of KTX stops outside city cores, KTX gains were at the expense of conventional rail, while intercity rail's modal share was little changed.
By 2007, provincial airports suffered from deficits after a drop in the number of passengers attributed to the KTX. With lower ticket prices, by 2008, KTX has swallowed up around half of the airlines' previous demand between Seoul and Busan (falling from 5.3 million passengers in 2003 to 2.4 million). Though some low-cost carriers failed and withdrew from the route, others still planned to enter competition even at the end of 2008. Budget airlines achieved a 5.6% growth in August 2009 over the same month a year earlier while KTX ridership decreased by 1.3%, a trend change credited to the opening of Seoul Subway Line 9, which improved Gimpo International Airport's connection to southern Seoul.
In the first two months after the launch of the second phase of the Gyeongbu HSR, passenger numbers on flights between Gimpo and Ulsan Airports dropped 35.4% compared to the same period a year earlier, those between Gimpo and Pohang Airports 13.2%. Between Gimpo Airport and Busan's Gimhae International Airport, airline passenger numbers remained stable (+0.2%), as a consequence of a budget airline competing with large discounts and aggressive marketing. In the first month of Gyeongjeon KTX service, express bus services between Seoul and Masan or Changwon experienced 30–40% drops in ridership.
Technical and operational issues
State of infrastructure
Lawmakers criticised the safety of Korail's tunnels after the Ministry of Construction and Transportation submitted data to the National Assembly on June 13, 2005. The ministry added fire prevention standards to high-speed line design standards only in November 2003, thus they weren't applied to the by then finished tunnels of the first phase of KTX. Consequently, few tunnels had emergency exits, and in high-speed railway tunnels, the average walking distance in case of an emergency was , with a maximum of , against a norm of emergency exits every in other countries. A contingency plan for fires in KTX tunnels was incorporated into a national disaster manual in November 2005.
On October 5, 2008, it was revealed by lawmakers that inside Hwanghak Tunnel, from December 2004, inspectors have monitored the progression of several cracks and minor track displacements, which continued after maintenance work in March–April 2007 and again in March 2008. The operator claimed that a February 2007 on-site inspection found the problems not safety-relevant, but pledged further maintenance, and an investigation into the causes was launched. Tunnel reinforcement was under way in 2010.
Incidents and accidents
Operation irregularities mostly concerned the rolling stock, but also signalling, power glitches and track problems. The number of incidents decreased from 28 in the first month to 8 in the fifth. The failure rate decreased sharply by the fifth year of operation. Later, in the first eight months of regular service until October 2010, KTX-II trains broke down 12 times. Causes for breakdowns in the first years of operation involved inexperienced staff and insufficient inspection during maintenance.
Lawmakers from the Grand National Party published an investigation in October 2006 and expressed concern about the practice to use parts from other trains for spare parts, but Korail stated that that is standard practice in case of urgency with no safety effect, and the supply of spare parts is secured. Korail is also conducting a localisation program to develop replacements for two dozen imported parts.
On June 13, 2007, near Cheongdo on the upgraded Daegu–Busan section, a damper acting between two cars of a KTX train got free at one end due to a loose screw and hit the trackbed, throwing up ballast that hit cars and caused bruises to two people on the parallel road, until the train was stopped when passengers noticed smoke.
On November 3, 2007, an arriving KTX-I train collided with a parked KTX-I train inside Busan Station, resulting in material damage of 10 billion won and light injuries to two persons. The accident happened because the driver had fallen asleep and disabled the train protection system, and led to the trial and conviction of the driver. The railway union criticised single driver operation in conjunction with the two and a half hours rest time the driver had between shifts.
On February 11, 2011, a KTX-Sancheon train bound for Seoul from Busan derailed on a switch in a tunnel before Gwangmyeong Station, when travelling at around . No casualties were reported, only one passenger suffered slight injury, but KTX traffic was blocked for 29 hours while repairs were completed. Preliminary investigation indicated that the accident resulted from a series of human errors. Because workers improperly repaired a point along the tracks. Investigators found that the derailment was caused by a switch malfunction triggered by a loose nut from track, and suspected that a repairman failed to tighten it during maintenance the previous night. The switch's detectors signalled a problem earlier, however, a second maintenance crew failed to find the loose nut and didn't properly communicate the fact to the control center, which then allowed the train on the track. The rail union criticised Korail's use of hired repairmen. there were no problems with the train according to investigation.
On July 15, 2011, 150 passengers were evacuated from a train when smoke started coming out of the train when it arrived at Miryang station at 11:30 AM. On July 17, 2011, at around 11 AM, a train stopped abruptly and stranded some 400 passengers in the Hwanghak Tunnel for over an hour. The train resumed service after emergency repairs to a malfunctioning motor. A Korail spokesperson stated that the reason for the stop was due to "faults in the motor block that supplies power to the wheels". The same day, the air conditioning broke down on another train leaving Busan at 1:45 PM. Over 800 passengers were transferred to another train at Daejeon when the problem could not be fixed.
On December 7, 2018, a KTX train carrying 198 passengers derailed about five minutes after leaving Gangneung for Seoul injuring 15 passengers. The train was traveling at about 103 km/h when almost all of its cars left the rails.
On January 5, 2022, a KTX-Sancheon train bound for Busan from Seoul carrying 303 passengers and crew derailed at 12:58 PM while passing a tunnel in Yeongdong of North Chungcheong Province, about 215 kilometers south of Seoul, injuring 7 passengers. The train was traveling at about 200 km/h when it partially derailed, resulting in a bogie wheel from car number 4 running off the track before being violently ejected from the train, throwing up ballast and causing structural damage to train cars. Subsequent KTX traffic was rerouted via the standard line, resulting in severe delays. Initially, it was believed that the derailment was caused by the train colliding with debris while passing Yeongdong Tunnel. However, evidence gathered from further investigation show that the missing bogie wheel was found inside Otan Tunnel, which is about 4 km before Yeongdong Tunnel, leading the investigating team to believe the train derailed due to faults within the wheel bogie assembly rather than from impact with debris.
Passenger comfort and convenience
Passenger surveys in the first months found that the limited capacity of bus connections and the lack of subway connections for intermediate stations, especially the newly built stations Gwangmyeong and Cheonan-Asan, was the problem mentioned most often. A better connection to Cheonan-Asan Station was provided by an extension of Seoul Subway Line 1 along the Janghang Line, opened on December 14, 2008. Gwangmyeong Station was linked to the same subway line by a shuttle service on December 15, 2006, but it made little impact due to the longtime differences between KTX and subway train schedules.
The noise level in the trains during tunnel passages was also subject to passenger complaints. This was referred to as a tunnel effect; it referred to both noise and vibration of the train when traveling through two specific tunnels. The tunnel effect was specifically noted as a reason for passenger dissatisfaction. Sound waves that are generally dispersed in an open environment are reflected against the tunnel walls, which causes the sound waves to come in contact with the passenger cabin and produces noise.
A reduction by 3–4 dB was achieved by retrofitting all trains with longer mud flaps at car ends until May 2006 to smooth the airflow at the articulated car joints. However, measurements in 2009 found significantly higher interior noise levels at some locations in two tunnels. Window thickness and sound insulation was improved in the KTX-II. The rails for high-speed trains like the KTX are welded together via a special techniques that make the rail a solid continuous rail; this method reduces the noise volume, which is produced by the wheels' contact with the rail, but it is not fully eliminated.
The isolation of KTX-I trains against pressure variations during tunnel passages was insufficient for some passengers, leading to efforts to reinforce pressurization in newer generations of trains. Pressure variations have been known to cause passengers to experience ringing in their ears; the ventilation systems on the passenger cabins are sealed when the train enters a tunnel in order to reduce the pressure changes. Pressure variations were not the only train cabin-associated complaint; KTX passengers were also known to have been negatively affected by inconsistent speeds of the trains.
Some KTX passengers found high-speed travel in backwards facing seats dizzying. Along with dizziness, feelings of nausea, headache, and sleepiness could also be experienced. Motion sickness was also noted as having had a minimal effect on KTX passengers; however, it still made an impact on passenger ride comfort. When the original seats were selected for the KTX trains, the anthropometry of the main consumers, who were largely expected to be Korean, were not considered. The seat design was found to have a significant effect on how passengers on the KTX trains rated the experience of their trip. Among the various factors that were considered to be vectors of discomfort were the angle of joints and specific areas of pressure, which were discovered to be present after an analysis of questionnaires that were completed by recent passengers. The factors of the seats of concern to KTX passengers were the shape, pitch, width, and the amount of legroom between the rows of seats. Swivel seats, which can be turned into the direction of travel, installed only on First Class in KTX-I trains, were made standard on both classes on newer generations of trains.
Studies have shown that term "ride comfort" has been used as an all-encompassing term for the KTX passengers' over all experience on the trains. While the KTX train is based on the French TGV model, it is considered to be more comfortable. The passengers' overall experience with regards to over-all ride comfort has been looked at as a combination of their physical health and emotional state. Fares were not included in the aforementioned questionnaires on ride comfort as there were variations in pricing due to seat arrangement, as well as weekday/weekend rates.
| Technology | High-speed rail | null |
523680 | https://en.wikipedia.org/wiki/Stretcher | Stretcher | A stretcher, gurney, litter, or pram is an apparatus used for moving patients who require medical care. A basic type (cot or litter) must be carried by two or more people. A wheeled stretcher (known as a gurney, trolley, bed or cart) is often equipped with variable height frames, wheels, tracks, or skids.
Stretchers are primarily used in acute out-of-hospital care situations by emergency medical services (EMS), military, and search and rescue personnel. In medical forensics, the right arm of a corpse is left hanging off the stretcher to let paramedics know it is a deceased person. They are also used to hold prisoners during lethal injections in the United States.
History
An early stretcher, likely made of wicker over a frame, appears in a manuscript from . Simple stretchers were common with militaries right through the middle of the 20th century.
Gurney
Generally spelled gurney, but also guerney or girney. The first usage of the term for a wheeled stretcher is unclear, but it is believed to have been derived from Pacific Coast slang. Its use in a hospital context was established by the 1930s.
Classification
EMS stretchers used in ambulances have wheels that makes transportation over pavement easier, and have a lock inside the ambulance and straps to secure the patient during transport. An integral lug on the stretcher locks into a sprung latch within the ambulance in order to prevent movement during transport. Modern stretchers may also have battery-powered hydraulics to raise and collapse the legs automatically. This eases the workload on EMS personnel, who are statistically at high risk of back injury from repetitive raising and lowering of patients. Specialized bariatric stretchers are also available, which feature a wider frame and higher weight capacity for heavier patients. Stretchers are usually covered with a disposable sheet or wrapping, and are cleaned after each use to prevent the spread of infection. Shelves, hooks and poles for medical equipment and intravenous medication are also frequently included.
Standard stretchers have several adjustments. The bed can be raised or lowered to facilitate patient transfer. The head of the stretcher can be raised so that the patient is in a sitting position (especially important for those in respiratory distress) or lowered flat in order to perform CPR, or for patients with suspected spinal injury who must be transported on a spinal board. The feet can be raised to what is called the Trendelenburg position, indicated for patients in shock.
Some manufacturers have begun to offer hybrid devices that combine the functionality of a stretcher, a recliner chair, and a treatment or procedural table into one device.
Basic stretchers
Simple stretchers are the most rudimentary type. They are lightweight and portable, made of canvas or other synthetic material suspended between two poles or tubular aluminum frame. Many are stored as disaster supplies and are often former military equipment.
The folding stretcher, also known as a top deck or collapsible stretcher, is similar in design to the simple stretcher, but features one or more hinged points of articulation to allow the stretcher to be collapsed into a more compact form for easier handling or storage. Some models may even allow the patient to sit upright in a Fowler's or Semi-Fowler's position.
The Roberson orthopedic stretcher or scoop stretcher is used for lifting patients, for instance from the ground onto an ambulance stretcher or onto a spinal board. The two ends of the stretcher can be detached from each other, splitting the stretcher into two longitudinal halves. To load a patient, one or both ends of the stretcher are detached, the halves placed under the patient from either side and fastened back together. With obese patients, the possibility exists of accidentally pinching the patient's back when closing the stretcher, so care must be made not to injure them when carrying out this procedure.
The litter, also known as a basket stretcher or Stokes litter, is designed to be used where there are obstacles to movement or other hazards: for example, in confined spaces, on slopes, in wooded terrain. Typically it is shaped to accommodate an adult in a face up position and it is used in search and rescue operations. The person is strapped into the basket, making safe evacuation possible. The litter has raised sides and often includes a removable head/torso cover for patient protection. After the person is secured in the litter, the litter may be wheeled, carried by hand, mounted on an ATV, towed behind skis, snowmobile, or horse, lifted or lowered on high angle ropes, or hoisted by helicopter.
The WauK board is also designed for use in small spaces. The patient is secured to the board with straps. It has two wheels and a foldable footrest at one end, allowing the patient to be moved by one person, much as with a hand truck for moving cargo. It can also be used at a variety of angles, making it easier to traverse obstacles, such as tight stairwells.
The Neil Robertson stretcher is a stretcher designed for securely transporting injured individuals through narrow and confined spaces, such as steep ladders, small hatchways, and narrow passages. The Neil Robertson stretcher is widely used in maritime settings, particularly by naval forces and maritime rescue team.
Flexible stretchers
A flexible stretcher, also known by the brand names Reeves sleeve or SKED, is a stretcher that is often supported longitudinally by wooden or plastic planks. Essentially a tarpaulin with handles, it is primarily used to move a patient through confined spaces, e.g., a narrow hallway, or to lift obese patients. Reeves stretchers have six handholds, allowing multiple rescuers to assist extrication.
Wheeled stretchers
For ambulances, a collapsible wheeled stretcher, or gurney, is a type of stretcher on a variable-height wheeled frame. Normally, an integral lug on the stretcher locks into a sprung latch within the ambulance in order to prevent movement during transport, often referred to as antlers due to their shape. It is usually covered with a disposable sheet and cleaned after each patient in order to prevent the spread of infection. Its key value is to facilitate moving the patient and sheet onto a fixed bed or table on arrival at the emergency department. Both types may have straps to secure the patient.
Other types of stretchers
The Nimier stretcher (brancard Nimier) was a type of stretcher used by the French army during World War I. The casualty was placed on their back, but in a "seated position", (that is, the thighs were perpendicular to the abdomen). Thus, the stretcher was shorter and could turn in the trenches. This type of stretcher is rarely seen today.
| Technology | Equipment | null |
523770 | https://en.wikipedia.org/wiki/Caramelization | Caramelization | Caramelization is a process of browning of sugar used extensively in cooking for the resulting butter-like flavor and brown color. The brown colors are produced by three groups of polymers: (C24H36O18), (C36H50O25), and (C125H188O80). As the process occurs, volatile chemicals such as diacetyl (known for its intense butter-like taste) are released, producing the characteristic caramel flavor.
Like the Maillard reaction, caramelization is a type of non-enzymatic browning. Unlike the Maillard reaction, caramelization is pyrolytic, as opposed to being a reaction with amino acids.
When caramelization involves the disaccharide sucrose, it is broken down into the monosaccharides fructose and glucose.
Process
Caramelization is a complex, poorly understood process that produces hundreds of chemical products, and includes the following types of reactions:
equilibration of anomeric and ring forms
sucrose inversion to fructose and glucose
condensation reactions
intramolecular bonding
isomerization of aldoses to ketoses
dehydration reactions
fragmentation reactions
unsaturated polymer formation
Effects of caramelization
The process is temperature-dependent. Specific sugars each have their own point at which the reactions begin to proceed readily. Impurities in the sugar, such as the molasses remaining in brown sugar, greatly speed the reactions.
Caramelization reactions are also sensitive to the chemical environment, and the reaction rate, or temperature at which reactions occur most readily, can be altered by controlling the level of acidity (pH). The rate of caramelization is generally lowest at near-neutral acidity (pH around 7), and accelerated under both acidic (especially pH below 3) and basic (especially pH above 9) conditions.
Uses in food
Caramelization is used to produce several foods, including:
Caramel sauce, a sauce made with caramel
Confiture de lait and dulce de leche, caramelized, sweetened milk
Caramel candies
Crème caramel, and the similar crème brûlée, a custard dish topped with sugar caramelized with a blowtorch
Caramelized onions, which are used in dishes like French onion soup. Onions require 30 to 45 minutes of cooking to caramelize.
Caramelized potatoes
Caramelized pears
Cola, of which some brands use caramelized sugar in small amounts for color
Latik, a sweet syrup made of sugar and coconut milk which is used in several Filipino desserts.
Dodol, a type of toffee made with cane sugar, rice flour, and coconut milk originating from Indonesia.
| Physical sciences | Organic reactions | Chemistry |
523932 | https://en.wikipedia.org/wiki/Chinese%20Eastern%20Railway | Chinese Eastern Railway | The Chinese Eastern Railway or CER (, , or , Kitaysko-Vostochnaya Zheleznaya Doroga or KVZhD), is the historical name for a railway system in Northeast China (also known as Manchuria).
The Russian Empire constructed the line from 1897 to 1902 during the Great Game period. The Railway was a concession to Russia, and later the Soviet Union, granted by the Qing dynasty government of Imperial China. The system linked Chita with Vladivostok in the Russian Far East and with Port Arthur, then an Imperial Russian leased ice-free port. The T-shaped line consisted of three branches:
the western branch, now the Harbin–Manzhouli Railway
the eastern branch, now the Harbin–Suifenhe Railway
the southern branch, now part of the Beijing–Harbin Railway
which intersected in Harbin. Saint Petersburg administered the railway and the concession, known as the Chinese Eastern Railway Zone, from the city of Harbin, which grew into a major rail-hub.
The southern branch of the CER, known as the Japanese South Manchuria Railway from 1906, became a locus and partial casus belli for the Russo-Japanese War of 1904–1905, the 1929 Sino-Soviet Conflict, and the Second Sino-Japanese War of 1937–1945. The Soviet Union sold the railway to the Japanese puppet state of Manchukuo in 1935; later in 1945 the Soviets regained co-ownership of the railway by treaty. The Soviet Union returned the Chinese Eastern Railway to the People's Republic of China in 1952.
Name
The official Chinese name of this railway was Great Qing Eastern Provinces Railway (), also known as Eastern Qing Railway () or Eastern Provinces Railway(). After the Xinhai Revolution, the northern branches was renamed to Chinese Eastern Provinces Railway () in 1915, shortened form as ().
The southern branch was renamed to South Manchuria Railway (Japanese kyujitai/) after Japanese took over from Russians in 1905.
It is also known in English as the Chinese Far East Railway, Trans-Manchurian Railway and North Manchuria Railway.
History
The Chinese Eastern Railway, a single-track line, provided a shortcut for the world's longest railroad, the Trans-Siberian Railway, from near the Siberian city of Chita, across northern Manchuria via Harbin to the Russian port of Vladivostok. This route drastically reduced the travel distance required along the originally proposed main northern route to Vladivostok, which lay completely on Russian soil but was not completed until a decade after the Manchurian "shortcut".
In 1896 China granted a construction concession through northern Manchuria under the supervision of Vice Minister of Public Works Xu Jingcheng. Work on the CER began in July 1897 along the line Tarskaya (east of Chita) — Hailar — Harbin — Nikolsk-Ussuriski, and accelerated drastically after Russia concluded a 25-year lease of Liaodong from China in 1898. Officially, traffic on the line started in November 1901, but regular passenger traffic from St. Petersburg to Vladivostok across the Trans-Siberian railway did not commence until July 1903.
In 1898, construction of a 550-mile (880 km) spur line, most of which later formed the South Manchuria Railway, began at Harbin, leading southwards through Eastern Manchuria, along the Liaodong Peninsula, to the ice-free deep-water port at Lüshun, which Russia was fortifying and developing into a first-class strategic naval base and marine coaling station for its Far East Fleet and Merchant Marine. This town was known in the west as Port Arthur, and the Russo-Japanese War (1904–1905) was fought largely over who would possess this region and its excellent harbor, as well as whether it would remain open to traders of all nations (Open Door Policy).
The Chinese Eastern Railway was essentially completed in 1902, a few years earlier than the stretch around Lake Baikal. Until the Circumbaikal portion was completed (1904–1905; double-tracked, 1914), goods carried on the Trans-Siberian Railway had to be trans-shipped by ferry almost a hundred kilometers across the lake (from Port Baikal to Mysovaya).
The Chinese Eastern Railway became important in international relations. After the first Sino-Japanese War of 1894–1895, Russia gained the right to build the Chinese Eastern Railway in Manchuria. They had a large army and occupied Northern Manchuria, which was of some concern to the Japanese. Russia wanted the railway badly. It loaned money to China and promised to use the proposed railway to help defend China against Japan, in the secret Li–Lobanov Treaty of 1896. Construction started in 1898 and was completed in 1903.
In 1900 during the Boxer Rebellion – which was suppressed by the Eight Nation Alliance including Russia – Russia also launched a separate invasion of Manchuria sending 100,000 troops to protect their interests in the railroad. During the Russo-Japanese War, Russia lost both the Liaodong Peninsula and much of the South Manchurian branch to Japan. The rail line from Changchun to Lüshun — transferred to Japanese control — became the South Manchuria Railway.
During the Russian Civil War (1917–1924) the Russian part of the CER came under the administration of the White Army. From the 1919 Karakhan Manifesto to 1927, diplomats of the Soviet Union would promise to revoke concessions in China, but the Soviets secretly kept tsarist concessions such as the Chinese Eastern Railway, as well as consulates, barracks, and Orthodox churches. This led Chiang Kai-Shek — who pushed foreign powers such as Britain to return some of their concessions from 1925 to 1927 — to turn against his former Soviet ally in 1927, seizing Soviet legations. The Soviets would later fight an armed conflict to keep control over the northern CER in the Sino-Soviet conflict of 1929. while Japan maintained control of the southern spur line.
After the establishment of Manchukuo it was known as the North Manchuria Railway until 23 March 1935, when the USSR sold its rights to the railway to the Manchukuo government; it was then merged into the Manchukuo National Railway and converted to standard gauge in four hours on 31 August.
From August 1945, the CER again came under the joint control of the USSR and China. After World War II the Soviet government insisted on occupying the Liaodong Peninsula but allowed joint control over the Southern branch with China; all this together received the name of the "Chinese Changchun Railway" ().
In 1952, the Soviet Union transferred (free of charge) all of its rights to the Chinese Changchun Railway to the People's Republic of China.
Flags
The flag of the Chinese Eastern Railway is a combination of Chinese and Russian flags. It changed several times with the political changes of both owners. The first CER flag (1897–1915) was a combination of the triangular version of the flag of the Qing dynasty and the flag of Russia, with East Provinces Railway of Great Qing () in Chinese. The 1915–1925 flag replaced the flag of the Qing dynasty with a triangular version of the five-colored flag, with East Provinces Railway Company of China () in Chinese. The flag was changed again in 1925 and 1932, with the flag of the Soviet Union and the flag of Manchukuo added.
Fleet
The only train that covers the entire route is the train #19/20 "Vostok" (translated as "East") Moscow — Beijing. The trip from Moscow to Beijing takes 146 hours (6 days, 2 hours). The journey in the opposite direction lasts 143 hours (5 days, 23 hours).
There is also a train #653/654 Zabaikalsk — Manzhouli which one can use to cross the Russian-Chinese border. The trip takes 25 minutes.
| Technology | Railway lines | null |
524003 | https://en.wikipedia.org/wiki/Internal%20and%20external%20angles | Internal and external angles | In geometry, an angle of a polygon is formed by two adjacent sides. For a simple polygon (non-self-intersecting), regardless of whether it is convex or non-convex, this angle is called an internal angle (or interior angle) if a point within the angle is in the interior of the polygon. A polygon has exactly one internal angle per vertex.
If every internal angle of a simple polygon is less than a straight angle ( radians or 180°), then the polygon is called convex.
In contrast, an external angle (also called a turning angle or exterior angle) is an angle formed by one side of a simple polygon and a line extended from an adjacent side.
Properties
The sum of the internal angle and the external angle on the same vertex is π radians (180°).
The sum of all the internal angles of a simple polygon is π(n−2) radians or 180(n–2) degrees, where n is the number of sides. The formula can be proved by using mathematical induction: starting with a triangle, for which the angle sum is 180°, then replacing one side with two sides connected at another vertex, and so on.
The sum of the external angles of any simple polygon, if only one of the two external angles is assumed at each vertex, is 2π radians (360°).
The measure of the exterior angle at a vertex is unaffected by which side is extended: the two exterior angles that can be formed at a vertex by extending alternately one side or the other are vertical angles and thus are equal.
Extension to crossed polygons
The interior angle concept can be extended in a consistent way to crossed polygons such as star polygons by using the concept of directed angles. In general, the interior angle sum in degrees of any closed polygon, including crossed (self-intersecting) ones, is then given by 180(n–2k)°, where n is the number of vertices, and the strictly positive integer k is the number of total (360°) revolutions one undergoes by walking around the perimeter of the polygon. In other words, the sum of all the exterior angles is 2πk radians or 360k degrees. Example: for ordinary convex polygons and concave polygons, k = 1, since the exterior angle sum is 360°, and one undergoes only one full revolution by walking around the perimeter.
| Mathematics | Two-dimensional space | null |
524122 | https://en.wikipedia.org/wiki/Metre%E2%80%93tonne%E2%80%93second%20system%20of%20units | Metre–tonne–second system of units | The metre–tonne–second (MTS) system of units was invented in France (hence the derived unit names sthène and pièze) where it became the legal system between 1919 and 1961. It was adopted by the Soviet Union in 1933 and abolished there in 1955. It was a coherent metric system of units, much as SI (itself a refinement of the MKS system) and the centimetre–gram–second system (CGS), but with larger units for industrial use, whereas the CGS system was regarded as only really suitable for laboratory use.
Units
The base units of the MTS system are:
Length: metre
Mass: tonne,
1 t = 103 kg = 1 Mg
Time: second
Some common derived units:
Volume: cubic metre or stere
1 m3 ≡ 1 st
Force: sthène,
1 sn = 1 t⋅m/s2 = 103 N = 1 kN
Energy: sthène-metre = kilojoule,
1 sn⋅m = 1 t⋅m2/s2 = 103 J = 1 kJ
Power: sthène-metre per second = kilowatt,
1 sn⋅m/s = 1 t⋅m2/s3 = 103 W = 1 kW
Pressure: pièze,
1 pz = 1 t/m⋅s2 = 103 Pa = 1 kPa = 1 cbar (centibar)
| Physical sciences | Measurement systems | Basics and measurement |
524124 | https://en.wikipedia.org/wiki/Hypercharge | Hypercharge | In particle physics, the hypercharge (a portmanteau of hyperonic and charge) Y of a particle is a quantum number conserved under the strong interaction. The concept of hypercharge provides a single charge operator that accounts for properties of isospin, electric charge, and flavour. The hypercharge is useful to classify hadrons; the similarly named weak hypercharge has an analogous role in the electroweak interaction.
Definition
Hypercharge is one of two quantum numbers of the SU(3) model of hadrons, alongside isospin . The isospin alone was sufficient for two quark flavours — namely and — whereas presently 6 flavours of quarks are known.
SU(3) weight diagrams (see below) are 2 dimensional, with the coordinates referring to two quantum numbers: (also known as ), which is the component of isospin, and , which is the hypercharge (defined by strangeness , charm , bottomness , topness , and baryon number ). Mathematically, hypercharge is
Strong interactions conserve hypercharge (and weak hypercharge), but weak interactions do not.
Relation with electric charge and isospin
The Gell-Mann–Nishijima formula relates isospin and electric charge
where I3 is the third component of isospin and Q is the particle's charge.
Isospin creates multiplets of particles whose average charge is related to the hypercharge by:
since the hypercharge is the same for all members of a multiplet, and the average of the I3 values is 0.
These definitions in their original form hold only for the three lightest quarks.
SU(3) model in relation to hypercharge
The SU(2) model has multiplets characterized by a quantum number J, which is the total angular momentum. Each multiplet consists of substates with equally-spaced values of Jz, forming a symmetric arrangement seen in atomic spectra and isospin. This formalizes the observation that certain strong baryon decays were not observed, leading to the prediction of the mass, strangeness and charge of the baryon.
The SU(3) has supermultiplets containing SU(2) multiplets. SU(3) now needs two numbers to specify all its sub-states which are denoted by λ1 and λ2.
specifies the number of points in the topmost side of the hexagon while specifies the number of points on the bottom side.
Examples
The nucleon group (protons with and neutrons with ) have an average charge of , so they both have hypercharge (since baryon number and ). From the Gell-Mann–Nishijima formula we know that proton has isospin while neutron has
This also works for quarks: For the up quark, with a charge of , and an of , we deduce a hypercharge of , due to its baryon number (since three quarks make a baryon, each quark has a baryon number of ).
For a strange quark, with electric charge , a baryon number of , and strangeness −1, we get a hypercharge so we deduce that That means that a strange quark makes an isospin singlet of its own (the same happens with charm, bottom and top quarks), while up and down constitute an isospin doublet.
All other quarks have hypercharge .
Practical obsolescence
Hypercharge was a concept developed in the 1960s, to organize groups of particles in the "particle zoo" and to develop ad hoc conservation laws based on their observed transformations. With the advent of the quark model, it is now obvious that strong hypercharge, , is the following combination of the numbers of up (), down (), strange (), charm (), top () and bottom ():
In modern descriptions of hadron interaction, it has become more obvious to draw Feynman diagrams that trace through the individual constituent quarks (which are conserved) composing the interacting baryons and mesons, rather than bothering to count strong hypercharge quantum numbers. Weak hypercharge, however, remains an essential part of understanding the electroweak interaction.
| Physical sciences | Quantum numbers | Physics |
524283 | https://en.wikipedia.org/wiki/Merino | Merino | The Merino is a breed or group of breeds of domestic sheep, characterised by very fine soft wool. It was established in Spain near the end of the Middle Ages, and was for several centuries kept as a strict Spanish monopoly; exports of the breed were not allowed, and those who tried risked capital punishment. During the eighteenth century, flocks were sent to the courts of a number of European countries, including France (where they developed into the Rambouillet), Hungary, the Netherlands, Prussia, Saxony and Sweden.
The Merino subsequently spread to many parts of the world, including South Africa, Australia, and New Zealand. They are presently common in South Africa. Numerous recognised breeds, strains and variants have developed from the original type; these include, among others, the American Merino and Delaine Merino in the Americas, the Australian Merino, Booroola Merino and Peppin Merino in Oceania, and the Gentile di Puglia, Merinolandschaf and Rambouillet in Europe.
The Australian Poll Merino is a polled (hornless) variant. Rams of other Merino breeds have long, spiral horns which grow close to the head, while ewes are usually hornless.
History
Etymology
The name merino was not documented in Spain until the early 15th century, and its origin is disputed.
Two suggested origins for the Spanish word merino are given in:
It may be an adaptation to the sheep of the name of a Castilian official inspector (merino) over a merindad, who may have also inspected sheep pastures. This word is from the medieval Latin maiorinus, a steward or head official of a village, from maior, meaning "greater". However, there is no indication in any of the Leonese or Castilian law codes that this official, either named as maiorinus or merino had any duties connected with sheep, and the late date at which merino was first documented makes any connection with the name of an early medieval magistrate implausible.
It also may be from the name of an Imazighen tribe, the Marini (or in Spanish, Benimerines), who occupied parts of the southwest of the Iberian Peninsula during the 12th and 13th centuries. This view gains some support from the derivation of many medieval Spanish pastoral terms from Arabic or Berber languages. However, an etymology based on a 12th-century origin for Merino sheep when the Marinids were in Spain is unacceptable; the origin of the breed occurred much later.
Origin
The three theories of the origins of the Merino breed in Spain are: the importation of North African flocks in the 12th century; its origin and improvement in Extremadura in the 12th and 13th centuries; the selective crossbreeding of Spanish ewes with imported rams at several different periods, so that its characteristic fine wool was not fully developed until the 15th century or even later. The first theory accepts that the breed was improved by later importation of North African rams and the second accepts an initial stock of North African sheep related to types from Asia Minor, and both claim an early date and largely North African origin for the Merino breed.
Sheep were relatively unimportant in the Islamic Caliphate of Córdoba, and there is no record of extensive transhumance before the caliphate's fall in the 1030s. The Marinids, when a nomadic Zenata Berber tribe, held extensive sheep flocks in what is now Morocco, and its leaders who formed the Marinid Sultanate militarily intervened in southern Spain, supporting the Emirate of Granada several times in the late 13th and early 14th centuries. Although they may possibly have brought new breeds of sheep into Spain, there is no definite evidence that the Marinids did bring extensive flocks to Spain. As the Marinids arrived as an intervening military force, they were hardly in a position to protect extensive flocks and practice selective breeding.
The third theory, that the Merino breed was created in Spain over several centuries with a strong Spanish heritage, rather than simply being an existing North African strain that was imported in the 12th century, is supported both by recent genetic studies and the absence of definitive Merino wool before the 15th century. The predominant native sheep breed in Spain from pre-Roman times was the churro, a homogeneous group closely related to European sheep types north of the Pyrenees and bred mainly for meat and milk, with coarse, coloured wool. Churro wool had little value, except where its ewes had been crossed with a fine wool breed from southern Italy in Roman times. Genetic studies have shown that the Merino breed most probably developed by the crossing of churro ewes with a variety of rams of other breeds at different periods, including Italian rams in Roman times, North African rams in the medieval period, and English rams from fine-wool breeds in the 15th century.
Although Spain exported wool to England, the Low Countries and Italy in the 13th and 14th centuries, it was only used to make cheap cloth. The earliest evidence of fine Spanish wool exports was to Italy in the 1390s and Flanders in the 1420s, although in both cases fine English wool was preferred. Spain became noted for its fine wool (spinning count between 60s and 64s) in the late 15th century, and by the mid-16th century its Merino wool was acknowledged to equal that of the finest English wools.
The earliest documentary evidence for Merino wools in Italy dates to the 1400s, and in the 1420s and 1430s, Merino wools were being mixed with fine English wool in some towns in the Low Countries to produce high quality cloth. However, it was only in the mid-16th century that the most expensive grades of cloth could be made entirely from Merino wool, after its quality had improved to equal that of the finest English wools, which were in increasingly short supply at that time.
Preserved medieval woollen fabrics from the Low Countries show that, before the 16th century, only the best quality English wools had a fineness of staple comparable to modern Merino wool. The wide range of Spanish wools produced in the 13th and early 14th centuries were mostly used domestically for cheap, coarse and light fabrics, and were not Merino wools. Later in the 14th century, similar non-Merino wools were exported from the northern Castilian ports of San Sebastián, Santander, and Bilbao to England and the Low Countries to make coarse, cheap cloth. The quality of Spanish wools exported increased markedly in the late 15th century, as did their price, promoted by the efforts of the monarchs Ferdinand and Isabella to improve quality.
Spain built up a virtual monopoly in fine wool exports in the final decades of the 15th century and in the 16th century, creating a substantial source of income for Castile. In part, this was because most English wool was woven and made into textile goods within England by the 16th century, rather than being exported.
Many of the Castillian Merino flocks were owned by nobility or the church, although Alfonso X realised that granting the urban elites of the towns of Old Castile and León transhumant rights would create an additional source of royal income and counteract the power of the privileged orders. During the late 15th, 16th and early 17th century, two-thirds of the sheep migrating annually were held in flocks of less than 100 sheep and very few flocks exceeded 1,000 sheep. By the 18th century, there were fewer small owners, and several owners held flocks of more than 20,000 sheep, but owners of small to moderately-sized flocks remained, and the Mesta was never simply a combination of large owners.
The transhumant sheep grazed the southern Spanish plains in winter and the northern highlands in summer. The annual migrations to and from Castile and León, where the sheep were owned and where they had summer pasturage, was organised and controlled by the Mesta along designated sheep-walks, or cañadas reales and arranged for suitable grazing, water and rest stops in these routes, and for shearing when the flocks started their return north.
The three Merino strains that founded the world's Merino flocks are the Royal Escurial flocks, the Negretti and the Paula. Among Merino bloodlines stemming from Vermont in the US, three historical studs were highly important: Infantado, Montarcos and Aguires. In recent times, Merino and breeds deriving from Merino stocks have spread worldwide. However, there has been a substantial decline in the numbers of several European Merino breeds, which are now considered to be endangered breeds and are no longer the subject of genetic improvement. In Spain, there are now two populations, the commercial Merino flocks, most common in the province of Extremadura and an "historical" Spanish Merino strain, developed and conserved in a breeding centre near Cordoba. The commercial Merino flocks show considerable genetic diversity, probably because of their cross-breeding with non-Spanish Merino-derived breeds since the 1960s, to create a strain more suitable for meat production. The historical Spanish strain, bred from animals selected from the main traditional Spanish genetic lines to ensure the conservation of a purebred lineage, exhibits signs of inbreeding.
Before the 18th century, the export of Merinos from Spain was a crime punishable by death. In the 18th century, small exportation of Merinos from Spain and local sheep were used as the foundation of Merino flocks in other countries. In 1723, some were exported to Sweden, but the first major consignment of Escurials was sent by Charles III of Spain to his cousin, Prince Xavier the Elector of Saxony, in 1765. Further exportation of Escurials to Saxony occurred in 1774, to Hungary in 1775 and to Prussia in 1786. Later in 1786, Louis XVI of France received 366 sheep selected from 10 different cañadas; these founded the stud at the Royal Farm at Rambouillet. In addition to the fine wool breeds mentioned, other breeds derived from Merino stocks were developed to produce mutton, including the French Ile de France and Berrichon du Cher breeds. Merino sheep were also sent to Eastern Europe where their breeding began in Hungary in 1774
The Rambouillet stud enjoyed some undisclosed genetic development with some English long-wool genes contributing to the size and wool-type of the French sheep. Through one ram in particular named Emperor – imported to Australia in 1860 by the Peppin brothers of Wanganella, New South Wales – the Rambouillet stud had an enormous influence on the development of the Australian Merino.
Sir Joseph Banks procured two rams and four ewes in 1787 by way of Portugal, and in 1792 purchased 40 Negrettis for King George III to found the royal flock at Kew. In 1808, 2000 Paulas were imported.
The King of Spain also gave some Escurials to the Dutch government in 1790; these thrived in the Dutch Cape Colony (South Africa). In 1788, John MacArthur, from the Clan Arthur (or MacArthur Clan) introduced Merinos to Australia from South Africa.
From 1765, the Germans in Saxony crossed the Spanish Merino with the Saxon sheep to develop a dense, fine type of Merino (spinning count between 70s and 80s) adapted to its new environment. From 1778, the Saxon breeding center was operated in the Vorwerk Rennersdorf. It was administered from 1796 by Johann Gottfried Nake, who developed scientific crossing methods to further improve the Saxon Merino. By 1802, the region had four million Saxon Merino sheep, and was becoming the centre for stud Merino breeding, and German wool was considered to be the finest in the world.
In 1802, Colonel David Humphreys, United States Ambassador to Spain, introduced the Vermont strain into North America with an importation of 21 rams and 70 ewes from Portugal and a further importation of 100 Infantado Merinos in 1808. The British embargo on wool and wool clothing exports to the U.S. before the 1812 British/U.S. war led to a "Merino Craze", with William Jarvis of the Diplomatic Corps importing at least 3,500
sheep between 1809 and 1811 through Portugal.
The Napoleonic Wars (1793–1813) almost destroyed the Spanish Merino industry. The old cabañas or flocks were dispersed or slaughtered. From 1810 onwards, the Merino scene shifted to Germany, the United States and Australia. Saxony lifted the export ban on living Merinos after the Napoleonic wars. Highly decorated Saxon sheep breeder Nake from Rennersdorf had established a private sheep farm in Kleindrebnitz in 1811, but ironically after the success of his sheep export to Australia and Russia, failed with his own undertaking.
United States Merinos
Merino sheep were introduced to Vermont in 1812. This ultimately resulted in a boom-bust cycle for wool, which reached a price of 57 cents/pound in 1835. By 1837, 1,000,000 sheep were in the state. The price of wool dropped to 25 cents/pound in the late 1840s. The state could not withstand more efficient competition from the other states, and sheep-raising in Vermont collapsed. Many sheep farmers from Vermont migrated with their flocks to other parts of the United States.
Australian Merinos
Early history
About 70 native sheep, suitable only for mutton, survived the journey to Australia with the First Fleet, which arrived in late January 1788. A few months later, the flock had dwindled to just 28 ewes and one lamb.
In 1797, Governor King, Colonel Patterson, Captain Waterhouse and Kent purchased sheep in Cape Town from the widow of Colonel Gordon, commander of the Dutch garrison. When Waterhouse landed in Sydney, he sold his sheep to Captain John MacArthur, Samuel Marsden and Captain William Cox. Although the early origin of the Australian Merino breed involved different stocks from Cape Colony, England, Saxony, France and America and although different Merino strains are bred in Australia, the Australian Merino populations are genetically similar and distinct from all other Merino populations, indicating a common history after they arrived in Australia.
John and Elizabeth Macarthur
By 1810, Australia had 33,818 sheep. John MacArthur (who had been sent back from Australia to England following a duel with Colonel Patterson) brought seven rams and one ewe from the first dispersal sale of King George III stud in 1804. The next year, MacArthur and the sheep returned to Australia, Macarthur to reunite with his wife Elizabeth, who had been developing their flock in his absence. Macarthur is considered the father of the Australian Merino industry; in the long term, however, his sheep had very little influence on the development of the Australian Merino.
Macarthur pioneered the introduction of Saxon Merinos with importation from the Electoral flock in 1812. The first Australian wool boom occurred in 1813, when the Great Dividing Range was crossed. During the 1820s, interest in Merino sheep increased. MacArthur showed and sold 39 rams in October 1820, grossing £510/16/5. In 1823, at the first sheep show held in Australia, a gold medal was awarded to W. Riley ('Raby') for importing the most Saxons; W. Riley also imported cashmere goats into Australia.
Eliza and John Furlong
Two of Eliza Furlong's (sometimes spelt Forlong or Forlonge) children had died from consumption, and she was determined to protect her surviving two sons by living in a warm climate and finding them outdoor occupations. Her husband John, a Scottish businessman, had noticed wool from the Electorate of Saxony sold for much higher prices than wools from New South Wales. The family decided on sheep farming in Australia for their new business. In 1826, Eliza walked over through villages in Saxony and Prussia, selecting fine Saxon Merino sheep. Her sons, Andrew and William, studied sheep breeding and wool classing. The selected 100 sheep were driven (herded) to Hamburg and shipped to Hull. Thence, Eliza and her two sons walked them to Scotland for shipment to Australia. In Scotland, the new Australia Company, which was established in Britain, bought the first shipment, so Eliza repeated the journey twice more. Each time, she gathered a flock for her sons. The sons were sent to New South Wales, but were persuaded to stop in Tasmania with the sheep, where Eliza and her husband joined them.
The Age in 1908 described Eliza Furlong as someone who had 'notably stimulated and largely helped to mould the prosperity of an entire state and her name deserved to live for all time in our history' (reprinted Wagga Wagga Daily Advertiser 27 January 1989).
John Murray
There were nearly 2 million sheep in Australia by 1830, and by 1836, Australia had won the wool trade war with Germany, mainly because of Germany's preoccupation with fineness. German manufacturers commenced importing Australian wool in 1845. In 1841, at Mount Crawford in South Australia, Murray established a flock of Camden-blood ewes mated to Tasmanian rams. To broaden the wool and give the animals some size, it is thought some English Leicester blood was introduced. The resultant sheep were the foundation of many South Australian strong wool studs. His brother Alexander Borthwick Murray was also a highly successful breeder of Merino sheep.
The Peppin brothers
The Peppin brothers took a different approach to producing a hardier, longer-stapled, broader wool sheep. After purchasing Wanganella Station in the Riverina, they selected 200 station-bred ewes that thrived under local conditions and purchased 100 South Australian ewes bred at Cannally that were sired by an imported Rambouillet ram. The Peppin brothers mainly used Saxon and Rambouillet rams, importing four Rambouillet rams in 1860. One of these, Emperor, cut an 11.4 lb (5.1 kg clean) wool clip. They ran some Lincoln ewes, but their introduction into the flock is undocumented. In 1865, George Merriman founded the fine wool Merino Ravensworth Stud, part of which is the Merryville Stud at Yass, New South Wales. His son, Sir Walter Merriman, incorporated Peppin bloodlines into his breeding program.
Vermont Merinos in Australia
In the 1880s, Vermont rams were imported into Australia from the U.S.; since many Australian stud men believed these sheep would improve wool cuts, their use spread rapidly. Unfortunately, the fleece weight was high, but the clean yield low, the greater grease content increased the risk of fly strike, they had lower uneven wool quality, and lower lambing percentages. Their introduction had a devastating effect on many famous fine-wool studs.
In 1889, while Australian studs were being devastated by the imported Vermont rams, several U.S. Merino breeders formed the Rambouillet Association to prevent the destruction of the Rambouillet line in the U.S. , an estimated 50% of the sheep on the U.S. western ranges are of Rambouillet blood.
The federation drought (1901–1903) reduced the number of Australian sheep from 72 to 53 million and ended the Vermont era. The Peppin and Murray blood strain became dominant in the pastoral and wheat zones of Australia.
High price records
The world record price for a ram was A$450,000 for JC&S Lustre 53, which sold at the 1988 Merino ram sale at Adelaide, South Australia. In 2008, an Australian Merino ewe was sold for A$14,000 at the Sheep Show and auction held at Dubbo, New South Wales.
Events
The New England Merino Field Days, which display local studs, wool, and sheep, are held during January in even numbered years in and around the Walcha, New South Wales district. The Annual Wool Fashion Awards, which showcase the use of Merino wool by fashion designers, are hosted by the city of Armidale, New South Wales in March each year.
Animal welfare developments
In Australia, mulesing of Merino sheep is a common practice to reduce the incidence of flystrike. It has been attacked by animal rights and animal welfare activists, with PETA running a campaign against the practice in 2004. The PETA campaign targeted U.S. consumers by using graphic billboards in New York City. PETA threatened U.S. manufacturers with television advertisements showing their companies' support of mulesing. Fashion retailers including Abercrombie & Fitch Co., Gap Inc and Nordstrom and George (UK) stopped stocking Australian Merino wool products.
New Zealand banned mulesing on 1 October 2018.
Characteristics
Merino is an excellent forager and very adaptable. It is bred predominantly for its wool, and its carcass size is generally smaller than that of sheep bred for meat. South African Meat Merino (SAMM), American Rambouillet and German Merinofleischschaf have been bred to balance wool production and carcass quality.
Merino has been domesticated and bred in ways that would not allow them to survive well without regular shearing by their owners. They must be shorn at least once a year because their wool does not stop growing. If this is neglected, the overabundance of wool can cause heat stress, mobility issues, and even blindness.
Wool qualities
Merino wool is fine and soft. Staples are commonly long. A Saxon Merino produces of greasy wool a year, while a good quality Peppin Merino ram produces up to . Merino wool is generally less than 24 micron (μm) in diameter. Basic Merino types include: strong (broad) wool (23 - 24.5 μm), medium wool (21 - 22.9 μm), fine (18.6 - 20.9 μm), superfine (15 – 18.5 μm) and ultra-fine (11.5 - 15 μm).
| Biology and health sciences | Sheep | null |
524383 | https://en.wikipedia.org/wiki/Ranunculus | Ranunculus | Ranunculus is a large genus of about 1750 species of flowering plants in the family Ranunculaceae. Members of the genus are known as buttercups, spearworts and water crowfoots.
The genus is distributed worldwide, primarily in temperate and montane regions. The familiar and widespread buttercup of gardens throughout Northern Europe (and introduced elsewhere) is the creeping buttercup Ranunculus repens, which has extremely tough and tenacious roots. Two other species are also widespread, the bulbous buttercup Ranunculus bulbosus and the much taller meadow buttercup Ranunculus acris. In ornamental gardens, all three are often regarded as weeds.
Buttercups usually flower in the spring, but flowers may be found throughout the summer, especially where the plants are growing as opportunistic colonizers, as in the case of garden weeds.
The water crowfoots (Ranunculus subgenus Batrachium), which grow in still or running water, are sometimes treated in a separate genus Batrachium (from Greek , "frog"). They have two different leaf types, thread-like leaves underwater and broader floating leaves. In some species, such as R. aquatilis, a third, intermediate leaf type occurs.
Ranunculus species are used as food by the larvae of some Lepidoptera species including the Hebrew character and small angle shades. Some species are popular ornamental flowers in horticulture, with many cultivars selected for large and brightly coloured flowers.
Distribution
Buttercups are found in both hemispheres on all continents aside from Antarctica, and are primarily found in temperate or montane habitats. They likely originated in northern Eurasia during the late Eocene or Oligocene and rapidly radiated up to the present, dispersing worldwide. Fossil evidence suggests that despite no longer occurring there, they inhabited Antarctica up to the mid-late Pliocene, even while glaciations were rapidly altering the landscape.
Fossil record
Ranunculus gailensis and Ranunculus tanaiticus seed fossils have been described from the Pliocene Borsoni Formation in the Rhön Mountains, central Germany. Achenes labelled Ranunculus cf. tachiroei is known from the Pliocene of the Hengduan Mountains of China. Indeterminate achenes have been found from Neogene strata from the Meyer Desert Formation biota in the Transantarctic Mountains, which appear to have inhabited a periglacial environment. The oldest potential fossil is from the Late Eocene (initially identified as Miocene) Florissant Formation of Colorado, identified by Theodore Dru Alison Cockerell in 1922.
Description
Plant
Buttercups are mostly perennial, but occasionally annual or biennial, herbaceous, aquatic or terrestrial plants, often with leaves in a rosette at the base of the stem. In many perennial species runners are sent out that will develop new plants with roots and rosettes at the distanced nodes.
The leaves lack stipules, have petioles, are palmately veined, entire, more or less deeply incised, or compound, and leaflets or leaf segments may be very fine and linear in aquatic species.
Flowers
The hermaphrodite flowers are single or in a cyme, have usually five (but occasionally as few as three or as many as seven) sepals and usually, five yellow, greenish or white petals that are sometimes flushed with red, purple or pink (but the petals may be absent or have a different, sometimes much higher number).
At the base of each petal is usually one nectary gland that is naked or may be covered by a scale. Anthers may be few, but often many are arranged in a spiral, are yellow or sometimes white, and with yellow pollen. The sometimes few but mostly many green or yellow carpels are not fused and are also arranged in a spiral, mostly on a globe or dome-shaped receptacle.
Reflective petals
The petals of buttercups are often highly lustrous, especially in yellow species, owing to a special coloration mechanism: the petal's upper surface is very smooth causing a mirror-like reflection. The flash aids in attracting pollinating insects and temperature regulation of the flower's reproductive organs.
Fruit
The fruits (in this case called achenes) may be smooth or hairy, winged, nobby or have hooked spines.
Naming
The genus name Ranunculus is Late Latin for "little frog", the diminutive of rana. This probably refers to many species being found near water, like frogs.
The common name buttercup may derive from a false belief that the plants give butter its characteristic yellow hue (in fact it is poisonous to cows and other livestock). A popular children's game involves holding a buttercup up to the chin; a yellow reflection is supposed to indicate a fondness for butter. In ancient Rome, a species of buttercup was held to the skin by slaves attempting to remove forehead tattoos made by their owners.
In the interior of the Pacific Northwest of the United States, the buttercup is called "Coyote's eyes"— in Nez Perce and in Sahaptin. In the legend, Coyote was tossing his eyes up in the air and catching them again when Eagle snatched them. Unable to see, Coyote made eyes from the buttercup.
Splitting of the genus
Molecular investigation of the genus has revealed that Ranunculus is not monophyletic with respect to a number of other recognized genera in the family—e.g. Ceratocephala, Halerpestes, Hamadryas, Laccopetalum, Myosurus, Oxygraphis, Paroxygraphis and Trautvetteria. A proposal to split Ranunculus into several genera has thus been published in a 2010 classification for the tribe Ranunculeae. The split (and often re-recognized) genera include Arcteranthis Greene, Beckwithia Jeps., Callianthemoides Tamura, Coptidium (Prantl) Beurl. ex Rydb., Cyrtorhyncha Nutt. ex Torr. & A.Gray, Ficaria Guett., Krapfia DC., Kumlienia E. Greene and Peltocalathos Tamura. Not all taxonomists and users accept this splitting of the genus, and it can alternatively be treated in the broad sense.
Pharmacological activity
The most common uses of Ranunculus species in traditional medicines are as an antirheumatic, as a rubefacient, and to treat intermittent fever. The findings in some Ranunculus species of, for example, protoanemonin, anemonin, may justify the uses of these species against fever, rheumatism and rubefacient in Asian traditional medicines.
Toxicity
All Ranunculus (buttercup) species are poisonous when eaten fresh, but their acrid taste and the blistering of the mouth caused by their poison means they are usually left uneaten. Poisoning in livestock can occur where buttercups are abundant in overgrazed fields where little other edible plant growth is left, and the animals eat them out of desperation. Symptoms of poisoning include bloody diarrhea, excessive salivation, colic, and severe blistering of the mouth, mucous membranes and gastrointestinal tract. When Ranunculus plants are handled, naturally occurring ranunculin is broken down to form protoanemonin, which is known to cause contact dermatitis in humans and care should therefore be exercised in extensive handling of the plants. The toxins are degraded by drying, so hay containing dried buttercups is safe.
Species
| Biology and health sciences | Ranunculales | Plants |
524466 | https://en.wikipedia.org/wiki/Clique%20%28graph%20theory%29 | Clique (graph theory) | In graph theory, a clique ( or ) is a subset of vertices of an undirected graph such that every two distinct vertices in the clique are adjacent. That is, a clique of a graph is an induced subgraph of that is complete. Cliques are one of the basic concepts of graph theory and are used in many other mathematical problems and constructions on graphs. Cliques have also been studied in computer science: the task of finding whether there is a clique of a given size in a graph (the clique problem) is NP-complete, but despite this hardness result, many algorithms for finding cliques have been studied.
Although the study of complete subgraphs goes back at least to the graph-theoretic reformulation of Ramsey theory by , the term clique comes from , who used complete subgraphs in social networks to model cliques of people; that is, groups of people all of whom know each other. Cliques have many other applications in the sciences and particularly in bioinformatics.
Definitions
A clique, , in an undirected graph is a subset of the vertices, , such that every two distinct vertices are adjacent. This is equivalent to the condition that the induced subgraph of induced by is a complete graph. In some cases, the term clique may also refer to the subgraph directly.
A maximal clique is a clique that cannot be extended by including one more adjacent vertex, that is, a clique which does not exist exclusively within the vertex set of a larger clique. Some authors define cliques in a way that requires them to be maximal, and use other terminology for complete subgraphs that are not maximal.
A maximum clique of a graph, , is a clique, such that there is no clique with more vertices. Moreover, the clique number of a graph is the number of vertices in a maximum clique in .
The intersection number of is the smallest number of cliques that together cover all edges of .
The clique cover number of a graph is the smallest number of cliques of whose union covers the set of vertices of the graph.
A maximum clique transversal of a graph is a subset of vertices with the property that each maximum clique of the graph contains at least one vertex in the subset.
The opposite of a clique is an independent set, in the sense that every clique corresponds to an independent set in the complement graph. The clique cover problem concerns finding as few cliques as possible that include every vertex in the graph.
A related concept is a biclique, a complete bipartite subgraph. The bipartite dimension of a graph is the minimum number of bicliques needed to cover all the edges of the graph.
Mathematics
Mathematical results concerning cliques include the following.
Turán's theorem gives a lower bound on the size of a clique in dense graphs. If a graph has sufficiently many edges, it must contain a large clique. For instance, every graph with vertices and more than edges must contain a three-vertex clique.
Ramsey's theorem states that every graph or its complement graph contains a clique with at least a logarithmic number of vertices.
According to a result of , a graph with 3n vertices can have at most 3n maximal cliques. The graphs meeting this bound are the Moon–Moser graphs K3,3,..., a special case of the Turán graphs arising as the extremal cases in Turán's theorem.
Hadwiger's conjecture, still unproven, relates the size of the largest clique minor in a graph (its Hadwiger number) to its chromatic number.
The Erdős–Faber–Lovász conjecture relates graph coloring to cliques.
The Erdős–Hajnal conjecture states that families of graphs defined by forbidden graph characterization have either large cliques or large cocliques.
Several important classes of graphs may be defined or characterized by their cliques:
A cluster graph is a graph whose connected components are cliques.
A block graph is a graph whose biconnected components are cliques.
A chordal graph is a graph whose vertices can be ordered into a perfect elimination ordering, an ordering such that the neighbors of each vertex v that come later than v in the ordering form a clique.
A cograph is a graph all of whose induced subgraphs have the property that any maximal clique intersects any maximal independent set in a single vertex.
An interval graph is a graph whose maximal cliques can be ordered in such a way that, for each vertex v, the cliques containing v are consecutive in the ordering.
A line graph is a graph whose edges can be covered by edge-disjoint cliques in such a way that each vertex belongs to exactly two of the cliques in the cover.
A perfect graph is a graph in which the clique number equals the chromatic number in every induced subgraph.
A split graph is a graph in which some clique contains at least one endpoint of every edge.
A triangle-free graph is a graph that has no cliques other than its vertices and edges.
Additionally, many other mathematical constructions involve cliques in graphs. Among them,
The clique complex of a graph G is an abstract simplicial complex X(G) with a simplex for every clique in G
A simplex graph is an undirected graph κ(G) with a vertex for every clique in a graph G and an edge connecting two cliques that differ by a single vertex. It is an example of median graph, and is associated with a median algebra on the cliques of a graph: the median m(A,B,C) of three cliques A, B, and C is the clique whose vertices belong to at least two of the cliques A, B, and C.
The clique-sum is a method for combining two graphs by merging them along a shared clique.
Clique-width is a notion of the complexity of a graph in terms of the minimum number of distinct vertex labels needed to build up the graph from disjoint unions, relabeling operations, and operations that connect all pairs of vertices with given labels. The graphs with clique-width one are exactly the disjoint unions of cliques.
The intersection number of a graph is the minimum number of cliques needed to cover all the graph's edges.
The clique graph of a graph is the intersection graph of its maximal cliques.
Closely related concepts to complete subgraphs are subdivisions of complete graphs and complete graph minors. In particular, Kuratowski's theorem and Wagner's theorem characterize planar graphs by forbidden complete and complete bipartite subdivisions and minors, respectively.
Computer science
In computer science, the clique problem is the computational problem of finding a maximum clique, or all cliques, in a given graph. It is NP-complete, one of Karp's 21 NP-complete problems. It is also fixed-parameter intractable, and hard to approximate. Nevertheless, many algorithms for computing cliques have been developed, either running in exponential time (such as the Bron–Kerbosch algorithm) or specialized to graph families such as planar graphs or perfect graphs for which the problem can be solved in polynomial time.
Applications
The word "clique", in its graph-theoretic usage, arose from the work of , who used complete subgraphs to model cliques (groups of people who all know each other) in social networks. The same definition was used by in an article using less technical terms. Both works deal with uncovering cliques in a social network using matrices. For continued efforts to model social cliques graph-theoretically, see e.g. , , and .
Many different problems from bioinformatics have been modeled using cliques. For instance, model the problem of clustering gene expression data as one of finding the minimum number of changes needed to transform a graph describing the data into a graph formed as the disjoint union of cliques; discuss a similar biclustering problem for expression data in which the clusters are required to be cliques. uses cliques to model ecological niches in food webs. describe the problem of inferring evolutionary trees as one of finding maximum cliques in a graph that has as its vertices characteristics of the species, where two vertices share an edge if there exists a perfect phylogeny combining those two characters. model protein structure prediction as a problem of finding cliques in a graph whose vertices represent positions of subunits of the protein. And by searching for cliques in a protein–protein interaction network, found clusters of proteins that interact closely with each other and have few interactions with proteins outside the cluster. Power graph analysis is a method for simplifying complex biological networks by finding cliques and related structures in these networks.
In electrical engineering, uses cliques to analyze communications networks, and use them to design efficient circuits for computing partially specified Boolean functions. Cliques have also been used in automatic test pattern generation: a large clique in an incompatibility graph of possible faults provides a lower bound on the size of a test set. describe an application of cliques in finding a hierarchical partition of an electronic circuit into smaller subunits.
In chemistry, use cliques to describe chemicals in a chemical database that have a high degree of similarity with a target structure. use cliques to model the positions in which two chemicals will bind to each other.
| Mathematics | Graph theory | null |
525073 | https://en.wikipedia.org/wiki/Bighorn%20sheep | Bighorn sheep | The bighorn sheep (Ovis canadensis) is a species of sheep native to North America. It is named for its large horns. A pair of horns may weigh up to ; the sheep typically weigh up to . Recent genetic testing indicates three distinct subspecies of Ovis canadensis, one of which is endangered: O. c. sierrae.
Sheep originally crossed to North America over the Bering Land Bridge from Siberia; the population in North America peaked in the millions, and the bighorn sheep entered into the mythology of Native Americans. By 1900, the population had crashed to several thousand due to diseases introduced through European livestock and overhunting.
Taxonomy and genetics
Ovis canadensis is one of two species of mountain sheep in North America; the other species being O. dalli, the Dall sheep. Wild sheep crossed the Bering land bridge from Siberia into Alaska during the Pleistocene (about 750,000 years ago); subsequently, they spread through western North America as far south as Baja California and northwestern mainland Mexico. Divergence from their closest Asian ancestor (snow sheep) occurred about 600,000 years ago. In North America, wild sheep diverged into two extant species — Dall sheep, which occupy Alaska and northwestern Canada, and bighorn sheep, which range from southwestern Canada to Mexico. However, the status of these species is questionable given that hybridization has occurred between them in their recent evolutionary history.
Former subspecies
In 1940, Ian McTaggart-Cowan split the species into seven subspecies, with the first three being mountain bighorns and the last four being desert bighorns:
Rocky Mountain bighorn sheep, O. c. canadensis, found from British Columbia to Arizona.
Badlands bighorn sheep (or Audubon's bighorn sheep), O. c. auduboni, occurred in North Dakota, South Dakota, Montana, Wyoming, and Nebraska. This subspecies has been extinct since 1925.
California bighorn sheep, O. c. californiana, found from British Columbia south to California and east to North Dakota. The definition of this subspecies has been updated (see below).
Desert bighorn sheep, O. c. nelsoni, the most common desert bighorn sheep, ranges from California through Arizona and in west Texas as the result of conservation and re-introduction efforts.
Mexican bighorn sheep, O. c. mexicana, ranges from Arizona and New Mexico south to Sonora and Chihuahua.
Peninsular bighorn sheep O. c. cremnobates, occur in the Peninsular Ranges of California and Baja California
Weems' bighorn sheep, O. c. weemsi, found in southern Baja California.
Current subspecies
Starting in 1993, Ramey and colleagues, using DNA testing, have shown this division into seven subspecies is largely illusory. Most scientists currently recognize three subspecies of bighorn. This taxonomy is supported by the most extensive genetics (microsatellite and mitochondrial DNA) study to date (2016) which found high divergence between Rocky Mountain and Sierra Nevada bighorn sheep, and that these two subspecies both diverged from desert bighorn before or during the Illinoian glaciation (about 315–94 thousand years ago). Thus, the three subspecies of O. canadensis are:
Rocky Mountain bighorn sheep (O. c. canadensis) – occupying the U.S. and Canadian Rocky Mountains, and the Northwestern United States.
Sierra Nevada bighorn sheep (O. c. sierrae) – formerly California bighorn sheep, a genetically distinct subspecies that only occurs in the Sierra Nevada in California. However, historic observer records suggest that bighorn sheep may have ranged as far west as the California Coastal Ranges, which are contiguous to the Sierra Nevada via the Transverse Ranges. An account of "wild sheep" in the vicinity of the Mission San Antonio near Jolon, California and the mountains around San Francisco Bay dates to circa 1769.
Desert bighorn sheep (O. c. nelsoni) – occurring throughout the desert regions of the Southwestern United States and Northwestern Mexico. The 2016 genetics study suggested a more modest divergence of this desert bighorn sheep into three lineages consistent with the earlier work of Cowan: Nelson's (O. c. nelsoni), Mexican (O. c. mexicana), and Peninsular (O. c. cremnobates). These three lineages occupy desert biomes that vary significantly in climate, suggesting exposure to different selection regimens.
In addition, two populations are currently considered endangered by the United States government:
Sierra Nevada bighorn sheep (O. c. sierrae),
Peninsular bighorn sheep, a distinct population segment of desert bighorn sheep (O. c. nelsoni)
Description
Bighorn sheep are named for the large, curved horns borne by the rams (males). Ewes (females) also have horns, but they are shorter and straighter. They range in color from light brown to grayish or dark, chocolate brown, with a white rump and lining on the backs of all four legs. Males typically weigh , are tall at the shoulder, and long from the nose to the tail. Females are typically , tall, and long. Male bighorn sheep have large horn cores, enlarged cornual and frontal sinuses, and internal bony septa. These adaptations serve to protect the brain by absorbing the impact of clashes. Bighorn sheep have preorbital glands on the anterior corner of each eye, inguinal glands in the groin, and pedal glands on each foot. Secretions from these glands may support dominance behaviors.
Bighorns from the Rocky Mountains are relatively large, with males that occasionally exceed and females that exceed . In contrast, Sierra Nevada bighorn males weigh up to only and females to . Males' horns can weigh up to , as much as all the bones in the male's body.
Natural history
Ecology
The Rocky Mountain and Sierra Nevada bighorn sheep occupy the cooler mountainous regions of Canada and the United States. In contrast, the desert bighorn sheep subspecies are indigenous to the hot desert ecosystems of the Southwestern United States and Mexico. Bighorn sheep inhabit alpine meadows, grassy mountain slopes, and foothill country near rugged, rocky cliffs and bluffs. Since bighorn sheep cannot move through deep snow, they prefer drier slopes, where the annual snowfall is less than about per year. A bighorn's winter range usually has lower elevations than its summer range.
Bighorn sheep are highly susceptible to certain diseases carried by domestic sheep, such as psoroptic scabies and pneumonia; additional mortality occurs as a result of accidents involving rock falls or falling off cliffs (a hazard of living in steep, rugged terrain). Bighorns are well adapted to climbing steep terrain, where they seek cover from predators. Predation primarily occurs with lambs, which are hunted by coyotes, bobcats, gray foxes, wolverines, jaguars, ocelots, lynxes, and golden eagles.
Bighorn sheep of all ages are threatened by black bears, grizzly bears, wolves, and especially mountain lions, which are perhaps best equipped with the agility to prey on them in uneven, rocky habitats. Fire suppression techniques may limit visibility through shrublands, and therefore increase cover and predation rates by mountain lions. Bighorn sheep are considered good indicators of land health because the species is sensitive to many human-induced environmental problems. In addition to their aesthetic value, bighorn sheep are considered desirable game animals by hunters.
Bighorn sheep graze on grasses and browse shrubs, particularly in fall and winter, and seek minerals at natural salt licks. Females tend to forage and walk, possibly to avoid predators and protect lambs, while males tend to eat and then rest and ruminate, which lends to more effective digestion and greater increase in body size.
Social structure and reproduction
Bighorn sheep live in large herds and do not typically follow a single leader ram, unlike the mouflon, the ancestor of the domestic sheep, which has a strict dominance hierarchy. Before the mating season or "rut", the rams attempt to establish a dominance hierarchy to determine access to ewes for mating. During the prerut period, most of the characteristic horn clashing occurs between rams, although this behavior may occur to a limited extent throughout the year. Bighorn sheep exhibit agonistic behavior: two competitors walk away from each other and then turn to face each other before jumping and lunging into headbutts. Rams' horns can frequently exhibit damage from repeated clashes. Females exhibit a stable, nonlinear hierarchy that correlates with age. Females may fight for high social status when they are integrated into the hierarchy at one to two years of age.
Rocky Mountain bighorn rams employ at least three different courting strategies. The most common and successful is the tending strategy, in which a ram follows and defends an estrous ewe. Tending takes considerable strength and vigilance, and ewes are most receptive to tending males, presumably feeling they are the most fit. Another tactic is coursing, when rams fight for an already tended ewe. Ewes typically avoid coursing males, so the strategy is ineffective. The rams also employ a blocking strategy. They prevent a ewe from accessing tending areas before she even enters estrus.
Bighorn ewes have a six-month gestation. In temperate climates, the peak of the rut occurs in November, with one, or rarely two, lambs being born in May. Most births occur in the first two weeks of the lambing period. Pregnant ewes of the Rocky Mountains migrate to alpine areas in spring, presumably to give birth in areas safer from predation, but are away from areas with good quality forage. Lambs born earlier in the season are more likely to survive than lambs born later. Lambs born late may not have access to sufficient milk, as their mothers are lactating at a time when food quality is lower. Newborn lambs weigh from and can walk within hours. The lambs are then weaned when they reach four to six months old. The lifespan of ewes is typically 10–14 years and 9–12 years for rams.
Infectious disease
Many bighorn sheep populations in the United States experience regular outbreaks of infectious pneumonia, which likely result from the introduction of bacterial pathogens (in particular, Mycoplasma ovipneumoniae, and some strains of Mannheimia haemolytica) carried asymptomatically in domestic sheep. Once introduced, pathogens can transmit rapidly through a bighorn population, resulting in all-age die-offs that sometimes kill up to 90% of the population. In the years following pathogen introduction, bighorn populations frequently experience multiple years of lamb pneumonia outbreaks. These outbreaks can severely limit recruitment and likely play a powerful role in slowing population growth.
Relationship with humans
Conservation
Bighorn sheep were widespread throughout the western United States, Canada, and northern Mexico two hundred years ago. The population was estimated to be 150,000 to 200,000. Unregulated hunting, habitat destruction, overgrazing of rangelands, and diseases contracted from domestic livestock all contributed to the decline, the most drastic occurring from about 1870 through 1950.
In 1936, the Arizona Boy Scouts mounted a statewide campaign to save the bighorn sheep. The scouts first became interested in the sheep through the efforts of Major Frederick Russell Burnham. Burnham observed that fewer than 150 of these sheep still lived in the Arizona mountains. The National Wildlife Federation, the Izaak Walton League, and the National Audubon Society also joined the effort. On January 18, 1939, over of land were set aside to create the Kofa National Wildlife Refuge and the Cabeza Prieta National Wildlife Refuge.
Many state and federal agencies have actively pursued the restoration of bighorn sheep since the 1940s. However, these efforts have met with limited success, and most of the historical range of bighorns remains unoccupied. Hunting for male bighorn sheep is allowed, but heavily regulated, in Canada and the United States.
In culture
Bighorn sheep were among the most admired animals of the Apsaalooka (Crow) people, and what is today called the Bighorn Mountain Range was central to the Apsaalooka tribal lands. In the Bighorn Canyon National Recreation Area book, storyteller Old Coyote describes a legend related to the bighorn sheep. A man possessed by evil spirits attempts to kill his heir by pushing the young man over a cliff, but the victim is saved by getting caught in trees. Rescued by bighorn sheep, the man takes the name of their leader, Big Metal. The other sheep grant him power, wisdom, sharp eyes, sure-footedness, keen ears, great strength, and a strong heart. Big Metal returns to his people with the message that the Apsaalooka people will survive only so long as the river winding out of the mountains is known as the Bighorn River.
Bighorn sheep are hunted for their meat and horns, used in ceremonies, as food, and as hunting trophies. They also serve as a source of ecotourism, as tourists come to see the bighorn sheep in their native habitat.
The Rocky Mountain bighorn sheep is the provincial mammal of Alberta and the state animal of Colorado and, as such, is incorporated into the symbol for the Colorado Division of Parks and Wildlife. The Desert bighorn sheep is the state mammal of Nevada.
The Bighorn sheep was featured in the children's book Buford the Little Bighorn (1967) by Bill Peet. The Bighorn sheep named Buford has a huge pair of horns in the Spring, Summer, Fall, and Winter, similar to Rudolph the Red-Nosed Reindeer.
Bighorn sheep were once known by the scientific identification "argali" or "argalia" due to assumption that they were the same animal as the Asiatic argali (Ovis ammon). Lewis and Clark recorded numerous sightings of O. canadensis in the journals of their exploration—sometimes using the name argalia. In addition, they recorded the use of bighorn sheep horns by the Shoshone in making composite bows. William Clark's Track Map produced after the expedition in 1814 indicated a tributary of the Yellowstone River named Argalia Creek and a tributary of the Missouri River named Argalia River, both in what is today Montana. Neither of these tributaries retained these names, however. The Bighorn River, another tributary of the Yellowstone, and its tributary stream, the Little Bighorn River, were both indicated on Clark's map and did retain their names, the latter being the namesake of the Battle of the Little Bighorn.
The Bighorn Ram was featured in a series of prints by artist Andy Warhol. In 1983, the artist was commissioned to create a portfolio of ten endangered species to raise environmental awareness. The portfolio, known as "Endangered Species" was created in support of the Endangered Species Act, which was passed by the U.S. Congress in 1973. Other animals within the portfolio include the Siberian Tiger, Bald Eagle and the Giant Panda.
| Biology and health sciences | Bovidae | Animals |
525310 | https://en.wikipedia.org/wiki/Middle%20finger | Middle finger | The middle finger, long finger, second finger, third finger, toll finger or tall man is the third digit of the human hand, located between the index finger and the ring finger. It is typically the longest digit. In anatomy, it is also called the third finger, digitus medius, digitus tertius or digitus III.
Overview
In Western countries, extending the middle finger (either by itself, or along with the index finger in the United Kingdom: see V sign) is an offensive and obscene gesture, widely recognized as a form of insult, due to its resemblance of an erect penis. It is known, colloquially, as "flipping the bird", "flipping (someone) off", or "giving (someone) the finger".
The middle finger is often used for finger snapping together with the thumb.
| Biology and health sciences | Human anatomy | Health |
525465 | https://en.wikipedia.org/wiki/Lock%20and%20key | Lock and key | A lock is a mechanical or electronic fastening device that is released by a physical object (such as a key, keycard, fingerprint, RFID card, security token or coin), by supplying secret information (such as a number or letter permutation or password), by a combination thereof, or it may only be able to be opened from one side, such as a door chain.
A key is a device that is used to operate a lock (to lock or unlock it). A typical key is a small piece of metal consisting of two parts: the bit or blade, which slides into the keyway of the lock and distinguishes between different keys, and the bow, which is left protruding so that torque can be applied by the user. In its simplest implementation, a key operates one lock or set of locks that are keyed alike, a lock/key system where each similarly keyed lock requires the same, unique key.
The key serves as a security token for access to the locked area; locks are meant to only allow persons having the correct key to open it and gain access. In more complex mechanical lock/key systems, two different keys, one of which is known as the master key, serve to open the lock. Common metals include brass, plated brass, nickel silver, and steel. The act of opening a lock without a key is called lock picking.
History
Premodern history
Locks have been in use for over 6000 years, with one early example discovered in the ruins of Nineveh, the capital of ancient Assyria. Locks such as this were developed into the Egyptian wooden pin lock, which consisted of a bolt, door fixture or attachment, and key. When the key was inserted, pins within the fixture were lifted out of drilled holes within the bolt, allowing it to move. When the key was removed, the pins fell part-way into the bolt, preventing movement.
The warded lock was also present from antiquity and remains the most recognizable lock and key design in the Western world. The first all-metal locks appeared between the years 870 and 900, and are attributed to English craftsmen. It is also said that the key was invented by Theodorus of Samos in the 6th century BC.
'The Romans invented metal locks and keys and the system of security provided by wards.'
Affluent Romans often kept their valuables in secure locked boxes within their households, and wore the keys as rings on their fingers. The practice had two benefits: It kept the key handy at all times, while signaling that the wearer was wealthy and important enough to have money and jewellery worth securing.
A special type of lock, dating back to the 17th–18th century, although potentially older as similar locks date back to the 14th century, can be found in the Beguinage of the Belgian city Lier. These locks are most likely Gothic locks, that were decorated with foliage, often in a V-shape surrounding the keyhole. They are often called drunk man's lock, as these locks were, according to certain sources, designed in such a way a person can still find the keyhole in the dark, although this might not be the case as the ornaments might have been purely aesthetic. In more recent times similar locks have been designed.
Modern locks
With the onset of the Industrial Revolution in the late 18th century and the concomitant development of precision engineering and component standardization, locks and keys were manufactured with increasing complexity and sophistication.
The lever tumbler lock, which uses a set of levers to prevent the bolt from moving in the lock, was invented by Robert Barron in 1778. His double acting lever lock required the lever to be lifted to a certain height by having a slot cut in the lever, so lifting the lever too far was as bad as not lifting the lever far enough. This type of lock is still used today.
The lever tumbler lock was greatly improved by Jeremiah Chubb in 1818. A burglary in Portsmouth Dockyard prompted the British Government to announce a competition to produce a lock that could be opened only with its own key. Chubb developed the Chubb detector lock, which incorporated an integral security feature that could frustrate unauthorized access attempts and would indicate to the lock's owner if it had been interfered with. Chubb was awarded £100 after a trained lock-picker failed to break the lock after 3 months.
In 1820, Jeremiah joined his brother Charles in starting their own lock company, Chubb. Chubb made various improvements to his lock: his 1824 improved design did not require a special regulator key to reset the lock; by 1847 his keys used six levers rather than four; and he later introduced a disc that allowed the key to pass but narrowed the field of view, hiding the levers from anybody attempting to pick the lock. The Chubb brothers also received a patent for the first burglar-resisting safe and began production in 1835.
The designs of Barron and Chubb were based on the use of movable levers, but Joseph Bramah, a prolific inventor, developed an alternative method in 1784. His lock used a cylindrical key with precise notches along the surface; these moved the metal slides that impeded the turning of the bolt into an exact alignment, allowing the lock to open. The lock was at the limits of the precision manufacturing capabilities of the time and was said by its inventor to be unpickable. In the same year Bramah started the Bramah Locks company at 124 Piccadilly, and displayed the "Challenge Lock" in the window of his shop from 1790, challenging "...the artist who can make an instrument that will pick or open this lock" for the reward of £200. The challenge stood for over 67 years until, at the Great Exhibition of 1851, the American locksmith Alfred Charles Hobbs was able to open the lock and, following some argument about the circumstances under which he had opened it, was awarded the prize. Hobbs' attempt required some 51 hours, spread over 16 days.
The earliest patent for a double-acting pin tumbler lock was granted to American physician Abraham O. Stansbury in England in 1805, but the modern version, still in use today, was invented by American Linus Yale Sr. in 1848. This lock design used pins of varying lengths to prevent the lock from opening without the correct key. In 1861, Linus Yale Jr. was inspired by the original 1840s pin-tumbler lock designed by his father, thus inventing and patenting a smaller flat key with serrated edges as well as pins of varying lengths within the lock itself, the same design of the pin-tumbler lock which still remains in use today. The modern Yale lock is essentially a more developed version of the Egyptian lock.
Despite some improvement in key design since, the majority of locks today are still variants of the designs invented by Bramah, Chubb and Yale.
Types of lock
Bicycle lock
Cam lock
Chamber lock
Child safety lock
Chubb detector lock
Combination lock
Cylinder lock
Dead bolt
Disc tumbler lock
Electric strike
Electromagnetic lock
Electronic lock
Lever tumbler lock
Luggage lock
Magnetic keyed lock
Mortise lock
Padlock
Pin tumbler lock
Police lock
Protector lock
Rim lock
Time lock
Warded lock
With physical keys
A warded lock uses a set of obstructions, or wards, to prevent the lock from opening unless the correct key is inserted. The key has notches or slots that correspond to the obstructions in the lock, allowing it to rotate freely inside the lock. Warded locks are typically reserved for low-security applications as a well-designed skeleton key can successfully open a wide variety of warded locks.
The pin tumbler lock uses a set of pins to prevent the lock from opening unless the correct key is inserted. The key has a series of grooves on either side of the key's blade that limit the type of lock the key can slide into. As the key slides into the lock, the horizontal grooves on the blade align with the wards in the keyway allowing or denying entry to the cylinder. A series of pointed teeth and notches on the blade, called bittings, then allow pins to move up and down until they are in line with the shear line of the inner and outer cylinder, allowing the cylinder or cam to rotate freely and the lock to open. An additional pin called the master pin is present between the key and driver pins in locks that accept master keys, to allow the plug to rotate at multiple pin elevations.
A wafer tumbler lock is similar to the pin tumbler lock and works on a similar principle. However, unlike the pin lock (where each pin consists of two or more pieces) each wafer is a single piece. The wafer tumbler lock is often incorrectly referred to as a disc tumbler lock, which uses an entirely different mechanism. The wafer lock is relatively inexpensive to produce and is often used in automobiles and cabinetry.
The disc tumbler lock or Abloy lock is composed of slotted rotating detainer discs.
The lever tumbler lock uses a set of levers to prevent the bolt from moving in the lock. In its simplest form, lifting the tumbler above a certain height will allow the bolt to slide past. Lever locks are commonly recessed inside wooden doors or on some older forms of padlocks, including fire brigade padlocks.
A magnetic keyed lock is a locking mechanism whereby the key utilizes magnets as part of the locking and unlocking mechanism. A magnetic key would use from one to many small magnets oriented so that the North and South poles would equate to a combination to push or pull the lock's internal tumblers thus releasing the lock.
With electronic keys
An electronic lock works by means of an electric current and is usually connected to an access control system. In addition to the pin and tumbler used in standard locks, electronic locks connect the bolt or cylinder to a motor within the door using a part called an actuator. Types of electronic locks include the following:
A keycard lock operates with a flat card of similar dimensions as a credit card. In order to open the door, one needs to successfully match the signature within the keycard.
The lock in a typical remote keyless system operates with a smart key radio transmitter. The lock typically accepts a particular valid code only once, and the smart key transmits a different rolling code every time the button is pressed.
Generally the car door can be opened with either a valid code by radio transmission, or with a (non-electronic) pin tumbler key.
The ignition switch may require a transponder car key to both open a pin tumbler lock and also transmit a valid code by radio transmission.
A smart lock is an electromechanics lock that gets instructions to lock and unlock the door from an authorized device using a cryptographic key and wireless protocol. Smart locks have begun to be used more commonly in residential areas, often controlled with smartphones. Smart locks are used in coworking spaces and offices to enable keyless office entry. In addition, electronic locks cannot be picked with conventional tools.
Locksmithing
Locksmithing is a traditional trade, and in most countries requires completion of an apprenticeship. The level of formal education required varies from country to country, from no qualifications required at all in the UK, to a simple training certificate awarded by an employer, to a full diploma from an engineering college. Locksmiths may be commercial (working out of a storefront), mobile (working out of a vehicle), institutional, or investigational (forensic locksmiths). They may specialize in one aspect of the skill, such as an automotive lock specialist, a master key system specialist or a safe technician. Many also act as security consultants, but not all security consultants have the skills and knowledge of a locksmith.
Historically, locksmiths constructed or repaired an entire lock, including its constituent parts. The rise of cheap mass production has made this less common; the vast majority of locks are repaired through like-for-like replacements, high-security safes and strongboxes being the most common exception. Many locksmiths also work on any existing door hardware, including door closers, hinges, electric strikes, and frame repairs, or service electronic locks by making keys for transponder-equipped vehicles and implementing access control systems.
Although the fitting and replacement of keys remains an important part of locksmithing, modern locksmiths are primarily involved in the installation of high quality lock-sets and the design, implementation, and management of keying and key control systems. Locksmiths are frequently required to determine the level of risk to an individual or institution and then recommend and implement appropriate combinations of equipment and policies to create a "security layer" that exceeds the reasonable gain of an intruder.
Key duplication
Traditional key cutting is the primary method of key duplication. It is a subtractive process named after the metalworking process of cutting, where a flat blank key is ground down to form the same shape as the template (original) key. The process roughly follows these stages:
The original key is fitted into a vise in a machine, with a blank attached to a parallel vise which is mechanically linked.
The original key is moved along a guide in a movement which follows the key's shape, while the blank is moved in the same pattern against a cutting wheel by the mechanical linkage between the vices.
After cutting, the new key is deburred by scrubbing it with a metal brush to remove particles of metal which could be dangerously sharp and foul locks.
Modern key cutting replaces the mechanical key following aspect with a process in which the original key is scanned electronically, processed by software, stored, then used to guide a cutting wheel when a key is produced. The capability to store electronic copies of the key's shape allows for key shapes to be stored for key cutting by any party that has access to the key image.
Different key cutting machines are more or less automated, using different milling or grinding equipment, and follow the design of early 20th century key duplicators.
Key duplication is available in many retail hardware stores and as a service of the specialized locksmith, though the correct key blank may not be available. More recently, online services for duplicating keys have become available.
Keyhole
A keyhole (or keyway) is a hole or aperture (as in a door or lock) for receiving a key. Lock keyway shapes vary widely with lock manufacturer, and many manufacturers have a number of unique profiles requiring a specifically milled key blank to engage the lock's tumblers.
Symbolism
Heraldry
Keys appear in various symbols and coats of arms, the best-known being that of the Holy See: derived from the phrase in Matthew 16:19 which promises Saint Peter, in Roman Catholic tradition the first pope, the Keys of Heaven. But this is by no means the only case.
Artwork
Some works of art associate keys with the Greek goddess of witchcraft, known as Hecate.
Palestinian key
The Palestinian key is the Palestinian collective symbol of their homes lost in the Nakba, when more than half of the population of Mandatory Palestine was expelled or fled violence in 1948 and were subsequently refused the right to return. Since 2016, a Palestinian restaurant in Doha, Qatar, holds the Guinness World Record for the world's largest key – 2.7 tonnes and 7.8 × 3 meters.
| Technology | Components_2 | null |
525538 | https://en.wikipedia.org/wiki/Topsoil | Topsoil | Topsoil is the upper layer of soil. It has the highest concentration of organic matter and microorganisms and is where most of the Earth's biological soil activity occurs.
Description
Topsoil is composed of mineral particles and organic matter and usually extends to a depth of 5-10 inches (13–25 cm). Together these make a substrate capable of holding water and air which encourages biological activity.
There are generally a high concentration of roots in topsoil since this is where plants obtain most of their vital nutrients. It also plays host to significant bacterial, fungal and entomological activity without which soil quality would degrade and become less suitable for plants. Bacteria and fungi can be essential in facilitating nutrient exchange with plants and in breaking down organic matter into a form that roots can absorb. Insects also play important roles in breaking down material and aerating and rotating the soil. Many species directly contribute to the health of the soil resulting in stronger plants. A healthy topsoil layer is a very rich microbiome that hosts a wide array of species.
Organic matter provides nutrition for living organisms and varies in quantity between different soils with the strength of the soil structure decreasing when more is present. It condenses and settles over time in different ways depending upon conditions such as beneath roadbeds and foundations vs uncovered and exposed to the elements. The structure becomes affected once the soil is dehydrated. Dehydrated topsoil volume substantially decreases and may suffer wind erosion.
Production
Topsoil is naturally produced in the process of soil formation or pedogenesis. Natural topsoil is mined and conditioned for human use and makes up the bulk of commercial topsoil available. The current rate of use and erosion outpaces soil generation.
It is possible to create artificial topsoil which supports some of the engineering or biological uses of topsoil. More traditional examples of artificial plant-growth media include terra preta and potting mix. Manufactured topsoil based on minerals, biosolids, compost and/or paper mill sludge is available commercially. A Victorian open-cut coal mine was rehabilitated with low-quality artificial topsoil made from local materials.
Classification
In soil classification systems, topsoil is known as the O Horizon or A Horizon. Soil horizons are layers parallel to the soil surface whose physical, chemical and biological characteristics differ from the layers above and beneath. The depth of the topsoil layer is measured as the depth of the surface to the first densely packed soil layer, known as subsoil.
In the United States, there is no federal, legal definition of the word topsoil when used in commerce.
Evaluation
Organisations such as the British Standards Institution (BSI) and the North Carolina Department of Agriculture publish guidelines for soil quality and the desired levels of topsoil nutrients broadly suitable for many plants.
Two common types of commercial topsoil are Bulk and Bagged Topsoil. The following table illustrates major differences between the two.
Alternatively the BSI relates the following values:
The preceding tables are for a multipurpose grade and certain levels can alter with regard to soil pH.
Standards also exist for specialist soils suitable for plants with specific needs including acidic or ericaceous soil and calcareous soil. These have different pH levels to typical soil and are meant for growing different plant species. Low fertility, low fertility acidic and low fertility calcareous are other soil classifications designed for plants which thrive in nutrient sparse soil.
Examples of specialist plants include the Venus flytrap which is found in low nitrogen and phosphorus environments so is less tolerant of highly nutrient rich environments than other plants and less able to compete in them. Whereas blueberries require ericaceous soil to grow well and clover grows well in calcareous soil. Soils must therefore be selected to suit the plants which are intended to be grown and hence standards are required.
Carbon to nitrogen ratio
Topsoil is the primary resource for plants to grow and crops to thrive. The main two parameters for this are carbon and nitrogen. The carbon provides energy and nitrogen is required for plants to build proteins and hence tissues. Plants require them in a range of ratios to enable suitable growth. An optimum figure for topsoil in the UK is a C:N ratio of less than 20:1. A sawdust base typically has a high C:N ratio in the order of 400:1 while an alfalfa hay has a low carbonaceous content and can typically have a C:N ratio around 12:1.
Commercial application
A variety of soil mixtures are sold commercially as topsoil. Typical uses for this product are improving gardens and lawns or for use in container gardens. Potting soil, compost, manure and peat are also sold for domestic uses with each having specific intended purposes. Topsoil products typically are not as suitable for potting plants or growing fruit and veg as potting soil or compost. Using it for this purpose can also work out prohibitively expensive compared to other alternatives.
Topsoil is also used for proper surface grading near residential buildings. In order to protect against flooding the International Residential Code requires a 2% slope () for the first ten feet away from the home.Energy Star requires a rate of .
Commercially available topsoil (manufactured or naturally occurring) in the United Kingdom must be classified to British Standard BS 3882, with the current version dated 2015. The standard has several classifications of topsoil with the final classification requiring material to meet certain threshold criteria such as nutrient content, extractable phytotoxic elements, particle size distribution, organic matter content, carbon:nitrogen ratio, electrical conductivity, loss on ignition, pH, chemical and physical contamination. The topsoil must be sampled in accordance with the British Standard and European Norm BS EN 12579:2013 Soil improvers and growing media – Sampling.
Erosion
Topsoil erosion occurs when the topsoil layer is blown or washed away. The estimated annual costs of public and environmental health losses related to soil erosion in the United States exceed $45 billion. Conventional industrial agriculture practices such as ploughing and spraying high quantities of synthetic liquid fertilisers can degrade the quality of the soil. Intensive farming methods to satisfy high food demands with high crop yields and growing crops in monocultures can deplete the soil nutrients and damage the soil microbiome. These factors can affect the consistency and quality of the soil resulting in increased erosion.
Surface runoff from farm fields is a type of nonpoint source pollution. Topsoil as well as farm fertilizers and other potential pollutants run off unprotected farm fields when heavy rains occur. This can result in polluting waterways and groundwater and may potentially contaminate drinking water sources. Algae blooms can occur when high quantities of nutrients flood rivers, lakes or oceans often as a result of farm runoff or from sewage. These harmful algal blooms can be toxic and have devastating impacts on ecosystems and wildlife. They are often referred to as red tides due to the presence of toxic red algae which can impact human food sources by contaminating seafood.
Sustainable techniques attempt to slow erosion through the use of cover crops in order to build organic matter in the soil. The United States loses almost 3 tons of topsoil per acre per year. of topsoil can take between 500 and 1,000 years to form naturally, making the rate of topsoil erosion a serious ecological concern. Based on 2014 trends, the world has about 60 years of topsoil left.
Conservation
| Physical sciences | Soil science | Earth science |
525737 | https://en.wikipedia.org/wiki/Berlin%20U-Bahn | Berlin U-Bahn | The Berlin U-Bahn (; short for , "underground railway") is a rapid transit system in Berlin, the capital and largest city of Germany, and a major part of the city's public transport system. Together with the S-Bahn, a network of suburban train lines, and a tram network that operates mostly in the eastern parts of the city, it serves as the main means of transport in the capital.
Opened in 1902, the serves 175 stations spread across nine lines, with a total track length of , about 80% of which is underground. Trains run every two to five minutes during peak hours, every five minutes for the rest of the day and every ten minutes in the evening. Over the course of a year, U-Bahn trains travel , and carry over 400 million passengers. In 2017, 553.1 million passengers rode the U-Bahn. The entire system is maintained and operated by the , commonly known as the BVG.
Designed to alleviate traffic flowing into and out of central Berlin, the U-Bahn was rapidly expanded until the city was divided into East and West Berlin at the end of World War II. Although the system remained open to residents of both sides at first, the construction of the Berlin Wall and the subsequent restrictions imposed by East Germany limited travel across the border. The East Berlin U-Bahn lines from West Berlin were severed, except for two West Berlin lines that ran through East Berlin (U6 and U8). These were allowed to pass through East Berlin without stopping at any of the stations, which were closed. Friedrichstraße was the exception because it was used as a transfer point between U6 and the West Berlin S-Bahn system, and a border crossing into East Berlin. The system was reopened completely following the fall of the Berlin Wall and German reunification.
The Berlin U-Bahn is the most extensive underground network in Germany.
In 2006, travel on the Berlin U-Bahn was equivalent to 122.2 million km (76 million mi) of car journeys.
History
The Berlin U-Bahn was built in three major phases:
Up to 1913: the construction of the (small profile) network in Berlin, Charlottenburg, Schöneberg, and Wilmersdorf;
Up to 1930: the introduction of the (large profile) network that established the first north–south lines;
From 1953 on: further development after World War II.
In a bid to secure its own improvement, Schöneberg also wanted a connection to Berlin. The elevated railway company did not believe such a line would be profitable, so the city built the first locally financed underground in Germany (intentionally using standard of Berlin U-Bahn rolling stock). It was opened on 1 December 1910. Just a few months earlier, work began on a fourth line to link Wilmersdorf in the southwest to the growing Berlin U-Bahn.
The early network ran mostly east to west, connecting the richer areas in and around Berlin, as these routes had been deemed the most profitable. In order to open up the network to more of the workers of Berlin, the city wanted north–south lines to be established. In 1920, the surrounding areas were annexed to form Groß-Berlin ("Greater Berlin Act"), removing the need for many negotiations, and giving the city much greater bargaining power over the private ("elevated railway company"). The city also mandated that new lines would use wider carriages—running on the same, standard-gauge track—to provide greater passenger capacity; these became known as the Großprofil ("large profile") network.
Construction of the ("North-South railway") connecting Wedding in the north to Tempelhof and Neukölln in the south had started in December 1912, but halted for the First World War. Work resumed in 1919, although the money shortage caused by hyperinflation slowed progress considerably. On 30 January 1923, the first section opened between Hallesches Tor and Stettiner Bahnhof (Naturkundemuseum), with a continuation to Seestraße following two months later. Desperately underfunded, the new line had to use trains from the old Kleinprofil network; the carriages exits had to be widened to fill the gap to the platforms with wooden boards that passengers jokingly referred to as Blumenbretter ("boards for flower pots"). The line branched at Belle-Alliance-Straße, now (Mehringdamm); the continuation south to Tempelhof opened on 22 December 1929, the branch to Grenzallee on 21 December 1930.
In 1912, plans were approved for AEG to build its own north–south underground line, named the after its termini, Gesundbrunnen and Neukölln, via Alexanderplatz. Financial difficulties stopped the construction in 1919; the liquidation of AEG-Schnellbahn-AG, and Berlin's commitment to the Nord-Süd-Bahn, prevented any further development until 1926. The first section opened on 17 July 1927 between Boddinstraße and Schönleinstraße, with the intermediate Hermannplatz becoming the first station at which passengers could transfer between two different Großprofil lines. The completed route was opened on 18 April 1930. Before control of the U-Bahn network was handed over completely to the BVG in 1929, the Hochbahngesellschaft started construction on a final line that, in contrast to its previous lines, was built as part of the Großprofil network. The major development was stopped in 1930.
The seizure of power by the National Socialists brought many changes that affected Germany, including the U-Bahn. Most notably, the new national flag was hung in every station, and two of the stations were renamed. Extensive plans—mostly the work of architect Albert Speer—were drawn up that included the construction of a circular line crossing the established U-Bahn lines, and new lines or extensions to many outlying districts. Despite such grand plans, no U-Bahn development occurred. In the Nazi period the only addition to Berlin's underground railways was North–South Tunnel of S-Bahn, opened 1936–1939.
During the Second World War, U-Bahn travel soared as car use fell, and many of the underground stations were used as air-raid shelters; however, Allied bombs damaged or destroyed large parts of the U-Bahn system. Although the damage was usually repaired fairly quickly, the reconstructions became more difficult as the war went on. Eventually, on 25 April 1945, the whole system ground to a halt when the power station supplying the network failed. Upon unconditional surrender of Nazi Germany following the Battle for Berlin there were 437 damaged points and 496 damaged vehicles.
The war had damaged or destroyed much of the network; however, of track and 93 stations were in use by the end of 1945, and the reconstruction was completed in 1950. Nevertheless, the consequent division of Berlin into East and West sectors brought further changes to the U-Bahn. Although the network spanned all sectors, and residents had freedom of movement, West Berliners increasingly avoided the Soviet sector and, from 1953, loudspeakers on the trains gave warnings when approaching the border, where passage of East Germans into the Western sectors also became subject to restrictions imposed by their government. There was a general strike on 17 June 1953 which closed the sections of the Berlin U-Bahn that traveled through East Berlin. Just after the strike, on the following day, train service on the line A was resumed and the service C was resumed to provide connections to Nordbahnhof and Friedrichstraße.
Between 1953 and 1955, the 200-Kilometre-Plan was drawn up, detailing the future development of the U-Bahn, which would grow to . Extending the C line to run from Tegel to Alt-Mariendorf was considered the highest priority: the northern extension to Tegel was opened on 31 May 1958. In order to circumvent East Berlin, and provide rapid-transport connections to the densely populated areas in Steglitz, Wedding, and Reinickendorf, a third north–south line was needed. The first section of line G was built between Leopoldplatz and Spichernstraße, with the intention of extending it at both ends. It had been planned to open the G line on 2 September 1961, but an earlier opening on 28 August was forced by the announcement of the construction of the Berlin Wall.
The next crisis was followed by the Berlin Wall construction on 13 August 1961, which had split the city between east and west. The U2 was split into two sections, and for the north–south lines, trains were not allowed to stop for passengers and become Geisterbahnhöfe ("ghost stations"), patrolled by armed East-German border guards. Only at Friedrichstraße, a designated border crossing point, were passengers allowed to disembark. A further consequence over the years is that most of the Berlin S-Bahn passengers boycotted the Deutsche Reichsbahn, and transferred to the U-Bahn with numerous expansion.
From 9 November 1989, following months of unrest, the travel restrictions placed upon East Germans were lifted. Tens of thousands of East Berliners heard the statement live on television and flooded the border checkpoints, demanding entry into West Berlin. Jannowitzbrücke, a former ghost station, was reopened two days later as an additional crossing point. It was the first station to be reopened after the opening of the Berlin Wall. Other stations, Rosenthaler Platz and Bernauer Straße on the U8 soon followed suit; and by 1 July 1990, all border controls were removed. In the decade following reunification, only three short extensions were made to U-Bahn lines.
In the 1990s some stations in the eastern portion of the city still sported bullet-riddled tiles at their entrances, a result of World War II battle damage during the Battle of Berlin. These were removed by 21 December 2004.
U-Bahn network
Routes
The U-Bahn has nine lines:
Stations
Among Berlin's 170 U-Bahn stations there are many with especially striking architecture or unusual design characteristics:
Hermannplatz station resembles something of a U-Bahn cathedral. The platform area is high, long and wide. It was built in connection with the construction of the first North-South Line (Nord-Süd-Bahn), now the U8. The architecturally important department store Karstadt adjacent to the station, was being constructed at the same time. Karstadt contributed a large sum of money towards the decoration of the station and was in return rewarded with direct access from the station to the store. Hermannplatz was also the first U-Bahn station in Berlin to be equipped with escalators. Today, Hermannplatz is a busy interchange between the U7 and U8.
Alexanderplatz station is another of the more notable U-Bahn stations in Berlin, and is an important interchange between three lines (U2, U5 and U8). The first part of the station was opened in 1913 along with an extension of today's U2 line. In the 1920s Alexanderplatz was completely redesigned, both above and below ground. The U-Bahn station was expanded to provide access to the new D (today's U8) and E (today's U5) lines, then under construction. The result was a station with a restrained blue-grey tiled colour-scheme and Berlin's first underground shopping facilities, designed by Alfred Grenander. Over the last few years Alexanderplatz station has, in stages, been restored; the work was due to be finished in 2007.
Wittenbergplatz station is also unusually designed. It opened in 1902 as a simple station with two side platforms, designed to plans created by Paul Wittig. The station was completely redesigned by Alfred Grenander in 1912, with five platform faces, accommodating two new lines, one to Dahlem on today's (U3), and the other to Kurfürstendamm, today's Uhlandstraße (Berlin U-Bahn) on the (U1). A provision for a sixth platform was included but has never been completed. The redesign also featured a new entrance building, which blended into the grand architectural styles of Wittenbergplatz and the nearby KaDeWe department store. The interior of the entrance building was again rebuilt after considerable war damage during World War II, this time in a contemporary 1950s style. This lasted until the early 1980s when the interior was retro-renovated back into its original style. Wittenbergplatz station was presented with a London style "Roundel type" station sign in 1952, the 50th Anniversary of the Berlin U-Bahn. Today's station is an interchange station between the U1, U2 and U3 lines.
The name of the Gleisdreieck (rail triangle) station is reminiscent of a construction which can only be imagined today. The wye was built in the opening year 1902. Plans for a redesign were made soon after, because the wye was already obsolete. An accident on 26 September 1908, which claimed 18 to 21 lives, was the final straw. The redesign and expansion of the transfer station, during which the station was still used, took until 1912. After World War II the station was put back into service on 21 October 1945 (lower platform) and 18 November 1945 (upper platform). However, service was interrupted again by the construction of the Berlin Wall. From 1972 onwards no trains ran on the lower platform, because servicing the U2 was no longer profitable due to the parallel traffic on the U1. The lower platform was reactivated in 1983, when the test line of the M-Bahn was built from the Gleisdreieck to the Kemperplatz station. It was broken down again after the fall of the Berlin Wall, since it obstructed parts of the reopened U2. Since 1993 the U1 and U2 trains both service the station again.
Tickets
Berlin public transit passes are available from many places, automated and non-automated, from BVG, Bahn, and authorized third-parties. The Ring-Bahn Line and the other S-Bahn lines are included, as are all U-Bahn lines, buses, trams, ferries, and most trains within the city limits: tickets are valid for all transportation considered part of the Berlin-Regional public transit system.
The Berlin U-Bahn mostly runs on an honor system and has been noted for its relative lack of turnstiles in its stations; instead transportation agents will inspect tickets and fine fare evaders.
Ride-passes (tickets) are available in fare classes: Adult and Reduced. Children between the ages of six and 14 and large dogs qualify for the reduced fare. Children below the age of six and small dogs travel free. There are senior discounts in the form of an annual ticket. Residents who have applied for and received a German Disability Identification card confirming 80% or more disability (ID's available from the Versorgungsamt, German Disability Office), can ride without a pass, including an additional person (as a helper). The disability identification card must be in the owner's possession when traveling.
With unemployment in the east averaging 15%, another common fare class in Berlin is the S(ocial)-Class. These identification cards are cleared through the normal government offices, then fulfilled at a BVG ride-pass non-automated location. Provided either by the Job Center (Arbeitsamt) for out-of-work residents or by the Sozialamt for people who cannot work or are disabled, the S-Class ride-passes normally restrict travel to the AB zones and must be renewed (a new pass purchased at a non-automated location) on the 1st of each month.
Additional passes are available for those which want to bring a bicycle on the public transit system. A bicycle-pass is included in the Student-class ride-pass, which is provided through the universities.
For small dogs which can be carried there is no additional fare requirement. For each "large dog", a reduced fare ride-pass must be purchased. Tourist ride-passes, all-day, group passes, and season passes include a dog fare.
BVG ride-passes are issued for specific periods of time, and most require validation with a stamping machine before they are first used. The validation shows the date and time of the first use, and where the ticket was validated (in code), and therefore when the ticket expires. For example, once validated, an all-day pass allows unlimited use from the time of purchase to 3:00 am the following day. Unlike most other metro systems, tickets in Berlin are not checked before entering tram, U-Bahn or S-Bahn stations. They are however checked by the bus drivers upon entering. On the tram, S-Bahn and U-Bahn, a proof-of-payment system is used: there are random spot checks inside by plain-clothed fare inspectors who have the right to demand to see each passenger's ticket. Passengers found without a ticket or an expired/invalid ticket are fined €60 per incident. The passenger may be required to pay on the spot, and is required on the spot to give a valid address to which the relevant fine notice can be mailed (it does not have to be in Germany). On the third incident, the BVG calls the offender to court, as there is now a history of 'riding without paying'.
Fare zones
Berlin is a part of the Verkehrsverbund Berlin-Brandenburg (Berlin-Brandenburg Transit Authority, VBB), which means ticketing and fare systems are unified with that of the surrounding state of Brandenburg. Berlin is divided into three fare zones, known as A, B, and C. Zone A is the area in the centre of Berlin and is demarcated by the S-Bahn urban rail ring line. Zone B covers the rest of the area within the city borders, and Zone C includes the immediate surroundings of Berlin. Zone C is divided into eight parts, each belonging to an administrative district. The Potsdam-Mittelmark area is included in the city district of Potsdam.
Tickets can be bought for specific fare zones, or multiple zones. Most passengers who live in Berlin buy AB fare zone tickets, while commuters coming in from the suburbs need ABC fare zone tickets. If a ticket not valid for travel in a tariff zone is checked by a ticket inspector, the passenger is subject to a fine.
Short-term tickets
Single-journey tickets (Einzeltickets) are issued for use within specific fare zones, namely AB, BC, and ABC. They are only valid for two hours after validation, and cannot be extended. The BVG also offers single-day tickets (Tageskarte), which are valid for the entire day when first validated until 3 a.m. the next morning.
Long-term tickets
Long-term paper tickets are issued with validity periods of seven days (7-Tage-Karte), one month (Monatskarten), or one year (Jahreskarte). The BVG is in the process of introducing the plastic MetroCard as a yearly ticket that also has additional features. The Metrocard also permits passengers to make reservations for hire cars at specific times, for example on weekends. It is expected that plastic Metrocards without such features will also be made available as they are more durable and ecofriendly than the paper tickets.
Tourist passes
The BVG offers tickets directed specifically for non-resident tourists of Berlin called the WelcomeCard and CityTourCard . WelcomeCards are valid for either 48 or 72 hours, and can be used by one adult and up to three children between the ages of six and 14. WelcomeCards are valid in fare zones ABC, and have the additional benefit of a reduction on entry fees to many museums and tourist attractions. See the Current Prices and Descriptions link for more information.
Underground facilities
Mobile phone network in 3G, 4G and 5G is in place throughout the entire U-Bahn networks. This system was in place by 1995 for the E-Plus network, and was one of the first metro systems in the world to allow mobile telephone use; by the late-1990s the other networks could be used in some portions as well. Since 2015, UMTS and LTE is also available for E-Plus and O2 (LTE since 2016) customers, and since 2020 mobile reception in some underground sections has also been extended to Deutsche Telekom and Vodafone Germany customers, with complete reception for the latter two telcos expected to be realised by mid-2021.
Many of the carriages on the U-Bahn feature small flat screen displays that feature news headlines from BZ, weekly weather forecasts, and ads for local businesses.
Most major interchange stations have large shopping concourses with banks, supermarkets, and fast food outlets.
Unused stations and tunnels
There are several stations, platforms and tunnels that were built in preparation for future U-Bahn extensions, and others that have been abandoned following planning changes. For example, platforms have already been provided for the planned "U3" at Potsdamer Platz on the planned line to Weißensee. It is unlikely that this line, which had the working title "U3" will ever be built, so the platforms have been partially converted into a location for events and exhibitions. The line number "U3" has been used to re-number the branch to Krumme Lanke, which had been part of "U1".
Line D, today's U8, was intended to run directly under Dresdner Straße via Oranienplatz to Kottbusser Tor. This segment of tunnel was abandoned in favour of a slightly less direct route in order to provide the former Wertheim department store at Moritzplatz with a direct connection. This involved the construction of a 90-degree curve of the line between Moritzplatz and Kottbusser Tor stations. The construction of the tunnel under Dresdner Straße had only been partially completed before abandonment, leaving it with only one track. This tunnel is separated into three parts, as it was blocked by a concrete wall where it crossed the border between East and West Berlin. Another concrete wall separates this tunnel, which now houses a transformer for an electricity supplier, from the never-completed Oranienplatz Station which is located partially under the square of the same name.
Stralauer Tor was a station on the eastern bank of the Spree between Warschauer Straße and Schlesisches Tor stations. It was completely destroyed in World War II. It had been opened in 1902 and was renamed Osthafen in 1924. Today, only struts on the viaduct remain to indicate its location. In the post-Second World War period it was not thought necessary to rebuild the station, due its close proximity to the Warschauer Straße station. Also its location was directly on the border between the Soviet and American sectors. Although a Berlin map dated 1946 shows the station renamed as Bersarinstraße after the Soviet General responsible for restoring civil administration of the city, this name was used later at another location.
Nürnberger Platz station was closed on 1 July 1959. It was replaced by two new stations on either side, Augsburger Straße and an interchange station to the U9 at Spichernstraße. Today, nothing remains of the station as a third track siding was constructed in its place.
Another tunnel, which once connected the U4 to its original depot and workshop at Otzenstraße (Schöneberg), is still in existence. The connection from Innsbrucker Platz station to the depot was severed when a deep level motorway underpass was constructed in the early 1970s; however, the continuation of the tunnel at Eisackstraße is still in existence for a distance of 270 metres and now ends at the former junction to the workshop of the Schöneberg line.
Platforms at five stations, Rathaus Steglitz, Schloßstraße, Walther-Schreiber-Platz, Innsbrucker Platz, and Kleistpark, were provided for the planned but never constructed U10. The U10 platform at Kleistpark has been converted into office space for the BVG. At Schloßstraße, U9 and U10 were planned to share two directional platforms at different levels; the would-be U10 tracks have been abandoned, leaving both platforms used by U9 trains only. The other U10 platforms remain unused and are not generally open to the public.
During the construction of Adenauerplatz (U7) station, which was built in conjunction with an underpass, platforms were also provided for a planned U1 extension from Uhlandstraße to Theodor-Heuss-Platz. A short tunnel section was also constructed in front of the Internationales Congress Centrum (ICC), beneath the Messedamm/Neue Kantstraße junction. This tunnel was built concurrently with a pedestrian subway and was also intended for the planned extension of the U1. The tunnel section, approximately long, ends at the location of the planned Messe station adjacent to Berlins central bus station (ZOB). The tunnel is used as a storage area for theater props.
At Jungfernheide station, double U-Bahn platforms similar to those at Schloßstraße were built for the planned extension of the U5. The unused platform sides are fenced off. The finished (U7) tunnel section which leads off towards Tegel airport is now used for firefighting exercises.
On 4 December 2020, the U5 extension between Alexanderplatz and Brandenburger Tor was opened. This included the new Unter den Linden station, which acts as a transfer point between the lines U5 and U6. Französische Straße station on the U6 was simultaneously closed due to its short distance to the new station.
Future development
Berlin's chronic financial problems make any expansion not mandated by the Hauptstadtvertrag—the document that regulates the necessary changes to the city as the capital of Germany—unlikely. Furthermore, there is still great rivalry for construction money between the U-Bahn and the S-Bahn. After the construction boom that followed the reunification of the city, enthusiasm for further growth has cooled off; many people feel that Berlin's needs are adequately met by the present U- and S-Bahn. As of 2020, the only proposals receiving serious consideration aim to facilitate travel around the existing system, such as moving Warschauer Straße's U-Bahn station closer to its S-Bahn station.
There are several long-term plans for the U-Bahn that have no estimated time of completion, most of which involve closing short gaps between stations, enabling them to connect to other lines. This would depend on demand, and new developments in the vicinity. New construction of U-Bahn lines is frequently the subject of political discussion with the Berlin chapters of the CDU, FDP and AfD who usually advocate in favor of U-Bahn expansion while the SPD, Alliance 90/The Greens and The Left typically advocate for tram construction instead.
After the last extension of U5 opened on 4 December 2020, there are no immediate plans to expand the metro system due to lack of budgetary conditions, although there are several extensions of railway lines that can be discussed over time:
Berlin Transport Minister Manja Schreiner (CDU) and Economy Minister Franziska Giffey (SPD) have underlined the local government's plans to extend the ends of each of the city's nine underground lines so that they reach the city's limits with the neighbouring state of Brandenburg.
"We must radically extend all the U-Bahn lines," Giffey told the Tagesspiegel newspaper. "We must offer Berliners a vision as to which routes we will tackle first," Schreiner added. Speaking to the dpa, Schreiner added that "Masterplan 2030" was crucial for many reasons: "More public transport means better climate protection, a better quality of life and more suitable mobility for everyone in the city."
Here's how the city plans to expand Berlin's nine existing underground lines - as well as building an additional two lines to serve travellers:
Ringlinie U0 - The outer Ringbahn
This is perhaps the biggest part of the project and one that will impact the most people. While Berlin's current Ringbahn - a circular line which rides an hour-long stretch around the city, connects U-Bahn lines to each other about mid-way through their routes - the local government plans for the U0 Ringlinie to connect the ends of each U-Bahn line that sprawls to near the outskirts of the city.
Since this part of the expansion project is particularly ambitious, it may be many, many years before you can step onto a U0 Ringlinie train.
U1 - Spandau to Weißensee
Currently, the U1 is simply a horizontal route which connects the east and west of central Berlin. Under the new plans, the U1 will reach Heerstraße in Spandau at one end and run through to Antonplatz in Weißensee at the other. With the plans, Antonplatz is set to become a new connection hotspot, where the U0 Ringlinie, the U1 and U3 will intersect.
U2 - Spandau to Pankow
This line will be expanded again into Spandau and towards the northeast to Pankow. The final stop in the northeast will be Pankow Kirche in the Pankow Altstadt.
U3 - Zehlendorf / Kleinmachnow to Falkenberg
This will be the only U-Bahn line that may even cross outside Berlin's borders with Brandenburg, reaching into Düppel-Kleinmachnow.
For now, though, the plans are just for the line to be extended to Mexikoplatz (S1), which reaches the border of Schlachtensee. If the funding is secured, the expected five-year-long construction process should begin swiftly and the U3 could reach Mexikoplatz by 2030.
U4 - Lichterfelde to Marzahn
This line will connect two very different parts of Berlin in what will be the biggest line extension of the project by far. At the moment the U4 is Berlin's baby U-Bahn line; the yellow one that stretches a modest four stops between Innsbrucker Platz and Nollendorfplatz without leaving the central southwest of the city.
Under the new plans, the line will be hugely extended at both ends, ultimately connecting Lichterfelde to Marzahn.
U5 - Charlottenburg to Hönow
Only recently was the U5 extended from Alexanderplatz to Hauptbahnhof and the line is already set for another development, but one not quite as ambitious as the U4 development.
Since the U5 already reaches quite far on its eastern side, to Hönow, it will only be extended in the west and even then only to Jungfernheide, the Ringbahn station that lies in northern Charlottenburg.
U6 - Tegel to Lichtenrade
As with the U5, the U6 will only be extended at one end of the line, in the south of the city. The new line will continue to run from Alt-Tegel, but instead of ending its journey at Alt-Mariendorf, will continue on to Naharlystraße in Lichtenrade.
U7 - Spandau to BER Airport
Already one of the city's longer U-Bahn routes, the U7 will be extended from Rathaus Spandau in the west, adding a new stop so it comes to meet the new U1 line at its terminal. In the southeast, the train will basically replace the current X7 bus route, which runs from Rudow to BER Airport.
U8 - Reinickendorf to Buckow
One of Berlin's most infamous lines, shamelessly voted to have the highest number of "disgusting" stations, the northern part of U8 will be expanded from Wittenau to reach the Märkisches Viertel in Reinickendorf. In the south, it will extend from Hermannstraße to Buckow-Süd.
U9 - Pankow to Buckow
Another of Berlin's most important north-south lines, the U9 will see considerable expansion at both ends. In the north the orange line will extend out from Wedding into Pankow and, in the south, the line will go quite far south of Steglitz to reach Buckower Chausee, south of Tempelhof-Schöneberg.
U10 - Alexanderplatz to Weißensee
And a new addition! The highly-awaited U10 should run from Alexanderplatz to Weißensee. Portions of the U10 have remained under the city's streets since the plans for the line were scrapped in the 1970s, earning it the name Phantomlinie (Phantom line). And when the U5 line was extended to Hauptbahnhof in 2020 an extra platform was built at U-Bahnhof Rotes Rathaus with the future U10 in mind.
New U-Bahn stops to expect on the U10 line are; Am Friedrichshain, Marienburger Straße, Danziger Straße, Greifswalder Straße, Gürtelstraße and Falkenberger Straße - but don't get too excited, the funds are yet to be secured for this one.
Rolling stock
The Berlin U-Bahn uses 750-volt DC electric trains that run on standard gauge ( ) tracks. The first trains were based on trams; they have a width of , and take their power from an upward facing third rail. To accommodate greater passenger numbers without lengthening the trains—which would require costly extended platforms—trains that ran on lines built after World War I were required to be wider. The original trains and lines, which continued to operate, were designated Kleinprofil (small profile), and the newer, wider trains and lines were designated Großprofil (large profile). Großprofil trains are wide, and take their power from a downward facing third rail. This is similar to New York City's A Division and B Division systems, where the B Division trains are wider than A Division trains (though B Division trains are also longer, while Großprofil trains are generally about the same length as Kleinprofil ones).
Although the two profiles are generally incompatible, Kleinprofil trains have been modified to run on Großprofil lines during three periods of economic difficulty. Between 1923 and 1927 on the Nord-Süd-Bahn, and between 1961 and 1978 on the E line, adapted Kleinprofil trains were used to compensate for the lack of new Großprofil trains: they were widened with wooden boards to reach the platforms; and had their power pickups adapted to accept power from the negatively charged downward-facing third rail, instead of positively charged upward-facing third rail. As of 2017, Class IK Kleinprofil trains are in operation on the Großprofil line U5, after refurbishing the existing F79 rolling stock was deemed unfeasible. They were widened with metal boards by on each side and elevated by to close the gap to the platforms; their power pickups were designed reversible to work on both profiles. As of October 2019, IK rolling stock is still used on the U5; it is intended to move the trains to Kleinprofil lines once new Großprofil rolling stock has been delivered.
As of 2007, Kleinprofil trains run on the U1, U2, U3, U4 and U5 lines; and Großprofil trains operate on the U5, U55, U6, U7, U8, and U9 routes.
Kleinprofil (small profile)
Kleinprofil trains are wide, and high. When the U-Bahn opened in 1902, forty-two multiple units, and twenty-one railroad cars, with a top speed of , had been built at the Warschauer Brücke workshop. In contrast to the earlier test vehicles, seating was placed along the walls, facing inward, which was considered more comfortable. Until 1927, U-Bahn trains had smoking compartments and third-class carriages. The trains were first updated in 1928; A-II carriages were distinguished by only having three windows, and two sliding doors.
After the division of the city, West Berlin upgraded its U-Bahn trains more rapidly than did East Berlin. The A3 type, introduced in 1960, was modelled on the Großprofil D type, and received regular modifications every few years. Meanwhile, A-I and A-II trains operated exclusively in East Berlin until 1975, when G-I trains, which had a top speed of , started to travel the Thälmannplatz–Pankow route. These were superseded in 1988 by the G-I/1 type, which used couplings that were incompatible with the older G-I carriages.
Following reunification, the A3L type was again upgraded as the A3L92. In 2000, prototypes for a Kleinprofil variant of the H series were built; the HK, the first Kleinprofil type to use AC induction motors like their large counterparts, differs from its Großprofil counterpart by not being fully interconnected—carriages are only interconnected within each of the two half-trains.
As of 2005, only trains of the HK, G-I/1 and A3(U/L) types are in active service.
From 2017, new IK-type trains will enter service to replace the remaining examples of type A3L71. Like HK-type trains they will be interconnected and as a result of their regenerative braking will recuperate up to 20% of the energy they require.
Großprofil (large profile)
Großprofil trains are wide, and high. The first sixteen multiple units and eight ordinary carriages entered active service on the Nord-Süd-Bahn in 1924, after a year of using modified Kleinprofil trains. Designated B-I, the cars were long and each had three sliding doors; the large elliptical windows at the front of the train earned them the nickname, Tunneleulen (tunnel owls). Upgraded B-II trains were introduced in 1927, and continued to be used until 1969. The first C-I trains were trialled in 1926, and two upgrades were produced before the end of the decade. The first U-Bahn trains to use aluminium in their construction, the C-IV types, were introduced in 1930. Many C-type trains were seized by Soviet forces in 1945, to be used in the Moscow Metro.
The first D-type trains, manufactured in 1957, were built from steel, making them very heavy and less efficient; however, the DL type that followed from 1965 used metals that were less dense, allowing a 26% reduction in weight. In East Berlin, D-type trains bought from the BVG were designated D-I. Difficulties there in trying to develop an E series of trains led, in 1962, to the conversion of S-Bahn type 168 trains for use on the E line. These E-III trains were desperately needed at the time to allow modified Kleinprofil trains to return to the increasingly busy A line but, following reunification, high running costs led to their retirement in 1994.
In West Berlin, the successor to the D-type was the F-type, which debuted in 1973. They varied from other models in having seats that were perpendicular to the sides of the train; from 1980, they also became the first U-Bahn trains to use three-phase electricity. In 1995, the original seating arrangement returned as the H series took up service. H-type trains are characterised by the interconnection of carriages throughout the length of the train; and they can only be removed from the tracks at main service depots.
As of 2005, only F, H, and a variation of the IK-type trains are in active service.
Depots
Depots of the Berlin U-Bahn fall into one of two classes: main workshops (, abbreviated as Hw); and service workshops (, abbreviated Bw). The main workshops are the only places where trains can be lifted from the tracks; they are used for the full inspections required every few years, and for any major work on trains. The service workshops only handle minor repairs and maintenance, such as replacing windows, or removing graffiti.
As of 2005, the only dedicated Kleinprofil depot is at Grunewald (Hw Gru/Bw Gru), which opened on 21 January 1913. The first Großprofil depot opened at Seestraße (Hw See/Bw See) in 1923, to service the Nord-Süd-Bahn. It has 17 tracks—2 for the main workshop, and 15 for the service workshop—but its inner-city location prevents any further expansion. Due to BVG budget cuts, the Seestraße depot also services Kleinprofil trains. Two further Großprofil service workshops are located at Friedrichsfelde (Bw Fri), and Britz-Süd (Bw Britz).
In the past, there were other workshops. The first opened in 1901 at Warschauer Brücke, and was the construction site for most of the early U-Bahn trains. The division of the U-Bahn network on 13 August 1961 forced its closure, although it was reopened in 1995 as a storage depot. A small depot operated at Krumme Lanke between 22 December 1929 and 1 May 1968; and, while the network was split, East Berlin's U-Bahn used the S-Bahn depot at Schöneweide, along with a small service workshop at Rosa-Luxemburg-Platz, which was closed following reunification.
Accidents
The Berlin ranks among the safest modes of transport: its history features few accidents.
The most severe accident occurred at the original (rail triangle), where the main and branch lines were connected by switches that allowed the tracks to cross. On 26 September 1908, a train driver missed a stop signal. As a result, two trains collided at the junction, and one fell off the viaduct. The accident killed eighteen people, and severely injured another twenty-one. 's triangular layout had already been deemed unsuitable for future developments; this incident—and a later, less-serious one—triggered its reconstruction as a multi-level station, starting in 1912.
On 30 June 1965, a train with brake failure stopped on the G line—today's U9—between and . Unaware of the faulty train, a mechanic working at the signal tower noticed that the signal for the affected section had been set to "Stop" for a long time. Thinking it was a fault of his, after several attempts he manually overrode the signal, in defiance of regulations that strictly prohibited such actions. The following train, which had been waiting at , then left the station on the same track. With emergency brakes unable to prevent the accident, the two trains collided. One passenger was killed in the crash, and 97 were injured. The mechanic was fined 600,000 DM.
Fires can be particularly dangerous and damaging within an underground system. In October 1972, two trains and a length of tunnel were completely destroyed when the trains caught fire; the reconstructed tunnel is clearly distinguishable from the old one. Another train burned out in the connecting tunnel between Klosterstraße and Alexanderplatz in 1987. On 8 July 2000, the last car of a GI/I train suffered a short circuit, burning out at the rear of the Deutsche Oper station. The single exit of the station was unreachable, forcing the passengers to run through the tunnel to reach the next emergency exit. The fire also damaged the station, which remained closed until that September. The Portuguese Ambassador, João Diogo Nunes Barata, presented the BVG with (tiled paintings), specially designed for the station, by the artist José de Guimarães. Installation of Portugal's gift to the city was completed on 30 October 2002.
As a consequence of the Deutsche Oper incident, BVG decided to post an employee at every station with only one exit until a second exit could be built. Over the following few years, many of those stations—including Britz-Süd, Schillingstraße, Viktoria-Luise-Platz, Uhlandstraße, and Theodor-Heuss-Platz—were retrofitted with additional exits. By June 2008, the only remaining stations with no second exit, Konstanzer Straße and Rudow, had been fitted with second exits. Despite these changes, several passenger organisations—such as Pro Bahn, and IGEB—demand that stations with exits in the middle of the platform also be fitted with additional emergency exits. Many stations are built this way; meeting those demands would place a heavy financial burden on both the BVG and the city.
The U6 saw a particularly costly, though casualty-free, incident on 25 March 2003. Scheduled repair work on the line limited the normal service to between Alt-Mariendorf and Kurt-Schumacher-Platz; one train then shuttled back and forth between Kurt-Schumacher-Platz and Holzhauser Straße, sharing a platform at Kurt-Schumacher-Platz with the normal-service trains departing for their return journey to Alt-Mariendorf. Needing to pass several stop signals on the shuttle service, the driver had been given special instructions how to proceed. Unfortunately, he ignored the signal at the entry to Kurt-Schumacher-Platz, and ploughed into the side of a train heading back to Alt-Mariendorf. The impact wrecked both trains, and caused considerable damage to the tracks. Normal service did not resume for two days, and the removal of the two wrecked trains—which, surprisingly, could still roll along the tracks—also took nearly 48 hours.
Films, music and merchandising
The Berlin U-Bahn has appeared in numerous films and music videos. Offering access to stations, tunnels, and trains, the BVG cooperates with film-makers, although a permit is required.
Whether set in Berlin or elsewhere, the U-Bahn has had at least a minor role in a large number of movies and television programmes, including Emil and the Detectives (2001), Otto – Der Film (1985), (1987) featuring Ingolf Lück, Run Lola Run (1998), and several Tatort episodes. The previously unused Reichstag station was used to shoot scenes of the movies Resident Evil and Equilibrium. The U-Bahn station Messe was used as coverage in the films Hanna and The Hunger Games: Mockingjay – Part 2.
Möbius 17, by Frank Esher Lämmer and Jo Preussler from Berlin, tells the story of an U-Bahn train that, caught in a Möbius strip, travels through alternate universes after a new line is built.
Alexanderplatz station plays an essential role in Berlin Alexanderplatz—a film of thirteen hour-long chapters and one epilogue—produced in 1980 by Rainer Werner Fassbinder, based on the book by Döblin. The film's scenes feature a recreation of the station as it was in 1928—rather darker and dirtier than in the 21st century. In the surrealistic two-hour epilogue, Fassbinder transforms parts of the station into a slaughterhouse where people are killed and dissected.
Since 2001, the Berlin U-Bahn has hosted the annual short-film festival Going Underground. Short films (up to 90 seconds long) are shown on the monitors found in many of the U-Bahn trains. Passengers on board vote for the festival winner.
Sandy Mölling, former singer of the pop band No Angels, shot the video for her single "Unnatural Blonde" in the U-Bahn station Deutsche Oper. Kate Ryan, Overground, Böhse Onkelz, Xavier Naidoo, Die Fantastischen Vier, and the DJ duo Blank & Jones have all used the U-Bahn and its stations for their videos as well.
"Linie 1", a musical performed by Berlin's Grips-Theater, is set completely in stations and trains of the Berlin U-Bahn; a movie version has also been produced.
In 2002, the BVG cooperated with design students in a project to create underwear with an U-Bahn theme, which, in English, they named "Underwear". They used the names of real stations that, in the context of underwear, appeared to be mild sexual double entendres: men's underpants bore labels with Rohrdamm (pipe dam), Onkel Toms Hütte (Uncle Tom's Cabin), and Krumme Lanke (crooked lake); the women's had Gleisdreieck (triangle track), and Jungfernheide (virgin heath). After the first series sold out quickly, several others were commissioned, such as Nothammer (emergency hammer), and Pendelverkehr (shuttle service; though Verkehr also means "intercourse" and Pendel also means "pendulum"). They were withdrawn from sale in 2004.
| Technology | Germany | null |
525887 | https://en.wikipedia.org/wiki/Density%20of%20states | Density of states | In condensed matter physics, the density of states (DOS) of a system describes the number of allowed modes or states per unit energy range. The density of states is defined as where is the number of states in the system of volume whose energies lie in the range from to . It is mathematically represented as a distribution by a probability density function, and it is generally an average over the space and time domains of the various states occupied by the system. The density of states is directly related to the dispersion relations of the properties of the system. High DOS at a specific energy level means that many states are available for occupation.
Generally, the density of states of matter is continuous. In isolated systems however, such as atoms or molecules in the gas phase, the density distribution is discrete, like a spectral density. Local variations, most often due to distortions of the original system, are often referred to as local densities of states (LDOSs).
Introduction
In quantum mechanical systems, waves, or wave-like particles, can occupy modes or states with wavelengths and propagation directions dictated by the system. For example, in some systems, the interatomic spacing and the atomic charge of a material might allow only electrons of certain wavelengths to exist. In other systems, the crystalline structure of a material might allow waves to propagate in one direction, while suppressing wave propagation in another direction. Often, only specific states are permitted. Thus, it can happen that many states are available for occupation at a specific energy level, while no states are available at other energy levels.
Looking at the density of states of electrons at the band edge between the valence and conduction bands in a semiconductor, for an electron in the conduction band, an increase of the electron energy makes more states available for occupation. Alternatively, the density of states is discontinuous for an interval of energy, which means that no states are available for electrons to occupy within the band gap of the material. This condition also means that an electron at the conduction band edge must lose at least the band gap energy of the material in order to transition to another state in the valence band.
This determines if the material is an insulator or a metal in the dimension of the propagation. The result of the number of states in a band is also useful for predicting the conduction properties. For example, in a one dimensional crystalline structure an odd number of electrons per atom results in a half-filled top band; there are free electrons at the Fermi level resulting in a metal. On the other hand, an even number of electrons exactly fills a whole number of bands, leaving the rest empty. If then the Fermi level lies in an occupied band gap between the highest occupied state and the lowest empty state, the material will be an insulator or semiconductor.
Depending on the quantum mechanical system, the density of states can be calculated for electrons, photons, or phonons, and can be given as a function of either energy or the wave vector . To convert between the DOS as a function of the energy and the DOS as a function of the wave vector, the system-specific energy dispersion relation between and must be known.
In general, the topological properties of the system such as the band structure, have a major impact on the properties of the density of states. The most well-known systems, like neutron matter in neutron stars and free electron gases in metals (examples of degenerate matter and a Fermi gas), have a 3-dimensional Euclidean topology. Less familiar systems, like two-dimensional electron gases (2DEG) in graphite layers and the quantum Hall effect system in MOSFET type devices, have a 2-dimensional Euclidean topology. Even less familiar are carbon nanotubes, the quantum wire and Luttinger liquid with their 1-dimensional topologies. Systems with 1D and 2D topologies are likely to become more common, assuming developments in nanotechnology and materials science proceed.
Definition
The density of states related to volume and countable energy levels is defined as:
Because the smallest allowed change of momentum for a particle in a box of dimension and length is , the volume-related density of states for continuous energy levels is obtained in the limit as
Here, is the spatial dimension of the considered system and the wave vector.
For isotropic one-dimensional systems with parabolic energy dispersion, the density of states is In two dimensions the density of states is a constant while in three dimensions it becomes
Equivalently, the density of states can also be understood as the derivative of the microcanonical partition function (that is, the total number of states with energy less than ) with respect to the energy:
The number of states with energy (degree of degeneracy) is given by:
where the last equality only applies when the mean value theorem for integrals is valid.
Symmetry
There is a large variety of systems and types of states for which DOS calculations can be done.
Some condensed matter systems possess a structural symmetry on the microscopic scale which can be exploited to simplify calculation of their densities of states. In spherically symmetric systems, the integrals of functions are one-dimensional because all variables in the calculation depend only on the radial parameter of the dispersion relation. Fluids, glasses and amorphous solids are examples of a symmetric system whose dispersion relations have a rotational symmetry.
Measurements on powders or polycrystalline samples require evaluation and calculation functions and integrals over the whole domain, most often a Brillouin zone, of the dispersion relations of the system of interest. Sometimes the symmetry of the system is high, which causes the shape of the functions describing the dispersion relations of the system to appear many times over the whole domain of the dispersion relation. In such cases the effort to calculate the DOS can be reduced by a great amount when the calculation is limited to a reduced zone or fundamental domain. The Brillouin zone of the face-centered cubic lattice (FCC) in the figure on the right has the 48-fold symmetry of the point group Oh with full octahedral symmetry. This configuration means that the integration over the whole domain of the Brillouin zone can be reduced to a 48-th part of the whole Brillouin zone. As a crystal structure periodic table shows, there are many elements with a FCC crystal structure, like diamond, silicon and platinum and their Brillouin zones and dispersion relations have this 48-fold symmetry. Two other familiar crystal structures are the body-centered cubic lattice (BCC) and hexagonal closed packed structures (HCP) with cubic and hexagonal lattices, respectively. The BCC structure has the 24-fold pyritohedral symmetry of the point group Th. The HCP structure has the 12-fold prismatic dihedral symmetry of the point group D3h. A complete list of symmetry properties of a point group can be found in point group character tables.
In general it is easier to calculate a DOS when the symmetry of the system is higher and the number of topological dimensions of the dispersion relation is lower. The DOS of dispersion relations with rotational symmetry can often be calculated analytically. This result is fortunate, since many materials of practical interest, such as steel and silicon, have high symmetry.
In anisotropic condensed matter systems such as a single crystal of a compound, the density of states could be different in one crystallographic direction than in another. These causes the anisotropic density of states to be more difficult to visualize, and might require methods such as calculating the DOS for particular points or directions only, or calculating the projected density of states (PDOS) to a particular crystal orientation.
k-space topologies
The density of states is dependent upon the dimensional limits of the object itself. In a system described by three orthogonal parameters (3 Dimension), the units of DOS is in a two dimensional system, the units of DOS is in a one dimensional system, the units of DOS is The referenced volume is the volume of -space; the space enclosed by the constant energy surface of the system derived through a dispersion relation that relates to . An example of a 3-dimensional -space is given in Fig. 1. It can be seen that the dimensionality of the system confines the momentum of particles inside the system.
Density of wave vector states (sphere)
The calculation for DOS starts by counting the allowed states at a certain that are contained within inside the volume of the system. This procedure is done by differentiating the whole k-space volume in n-dimensions at an arbitrary , with respect to . The volume, area or length in 3, 2 or 1-dimensional spherical -spaces are expressed by
for a -dimensional -space with the topologically determined constants
for linear, disk and spherical symmetrical shaped functions in 1, 2 and 3-dimensional Euclidean -spaces respectively.
According to this scheme, the density of wave vector states is, through differentiating with respect to , expressed by
The 1, 2 and 3-dimensional density of wave vector states for a line, disk, or sphere are explicitly written as
One state is large enough to contain particles having wavelength λ. The wavelength is related to through the relationship.
In a quantum system the length of λ will depend on a characteristic spacing of the system L that is confining the particles. Finally the density of states N is multiplied by a factor , where is a constant degeneracy factor that accounts for internal degrees of freedom due to such physical phenomena as spin or polarization. If no such phenomenon is present then . Vk is the volume in k-space whose wavevectors are smaller than the smallest possible wavevectors decided by the characteristic spacing of the system.
Density of energy states
To finish the calculation for DOS find the number of states per unit sample volume at an energy inside an interval . The general form of DOS of a system is given as
The scheme sketched so far only applies to monotonically rising and spherically symmetric dispersion relations. In general the dispersion relation is not spherically symmetric and in many cases it isn't continuously rising either. To express D as a function of E the inverse of the dispersion relation has to be substituted into the expression of as a function of k to get the expression of as a function of the energy. If the dispersion relation is not spherically symmetric or continuously rising and can't be inverted easily then in most cases the DOS has to be calculated numerically. More detailed derivations are available.
Dispersion relations
The dispersion relation for electrons in a solid is given by the electronic band structure.
The kinetic energy of a particle depends on the magnitude and direction of the wave vector k, the properties of the particle and the environment in which the particle is moving. For example, the kinetic energy of an electron in a Fermi gas is given by
where m is the electron mass. The dispersion relation is a spherically symmetric parabola and it is continuously rising so the DOS can be calculated easily.
For longitudinal phonons in a string of atoms the dispersion relation of the kinetic energy in a 1-dimensional k-space, as shown in Figure 2, is given by
where is the oscillator frequency, the mass of the atoms, the inter-atomic force constant and inter-atomic spacing. For small values of the dispersion relation is linear:
When the energy is
With the transformation and small this relation can be transformed to
Isotropic dispersion relations
The two examples mentioned here can be expressed like
This expression is a kind of dispersion relation because it interrelates two wave properties and it is isotropic because only the length and not the direction of the wave vector appears in the expression. The magnitude of the wave vector is related to the energy as:
Accordingly, the volume of n-dimensional -space containing wave vectors smaller than is:
Substitution of the isotropic energy relation gives the volume of occupied states
Differentiating this volume with respect to the energy gives an expression for the DOS of the isotropic dispersion relation
Parabolic dispersion
In the case of a parabolic dispersion relation (p = 2), such as applies to free electrons in a Fermi gas, the resulting density of states, , for electrons in a n-dimensional systems is
for , with for .
In 1-dimensional systems the DOS diverges at the bottom of the band as drops to . In 2-dimensional systems the DOS turns out to be independent of . Finally for 3-dimensional systems the DOS rises as the square root of the energy.
Including the prefactor , the expression for the 3D DOS is
where is the total volume, and includes the 2-fold spin degeneracy.
Linear dispersion
In the case of a linear relation (p = 1), such as applies to photons, acoustic phonons, or to some special kinds of electronic bands in a solid, the DOS in 1, 2 and 3 dimensional systems is related to the energy as:
Distribution functions
The density of states plays an important role in the kinetic theory of solids. The product of the density of states and the probability distribution function is the number of occupied states per unit volume at a given energy for a system in thermal equilibrium. This value is widely used to investigate various physical properties of matter. The following are examples, using two common distribution functions, of how applying a distribution function to the density of states can give rise to physical properties.
Fermi–Dirac statistics: The Fermi–Dirac probability distribution function, Fig. 4, is used to find the probability that a fermion occupies a specific quantum state in a system at thermal equilibrium. Fermions are particles which obey the Pauli exclusion principle (e.g. electrons, protons, neutrons). The distribution function can be written as
is the chemical potential (also denoted as EF and called the Fermi level when T=0), is the Boltzmann constant, and is temperature. Fig. 4 illustrates how the product of the Fermi-Dirac distribution function and the three-dimensional density of states for a semiconductor can give insight to physical properties such as carrier concentration and Energy band gaps.
Bose–Einstein statistics: The Bose–Einstein probability distribution function is used to find the probability that a boson occupies a specific quantum state in a system at thermal equilibrium. Bosons are particles which do not obey the Pauli exclusion principle (e.g. phonons and photons). The distribution function can be written as
From these two distributions it is possible to calculate properties such as the internal energy per unit volume , the number of particles , specific heat capacity , and thermal conductivity . The relationships between these properties and the product of the density of states and the probability distribution, denoting the density of states by instead of , are given by
is dimensionality, is sound velocity and is mean free path.
Applications
The density of states appears in many areas of physics, and helps to explain a number of quantum mechanical phenomena.
Quantization
Calculating the density of states for small structures shows that the distribution of electrons changes as dimensionality is reduced. For quantum wires, the DOS for certain energies actually becomes higher than the DOS for bulk semiconductors, and for quantum dots the electrons become quantized to certain energies.
Photonic crystals
The photon density of states can be manipulated by using periodic structures with length scales on the order of the wavelength of light. Some structures can completely inhibit the propagation of light of certain colors (energies), creating a photonic band gap: the DOS is zero for those photon energies. Other structures can inhibit the propagation of light only in certain directions to create mirrors, waveguides, and cavities. Such periodic structures are known as photonic crystals. In nanostructured media the concept of local density of states (LDOS) is often more relevant than that of DOS, as the DOS varies considerably from point to point.
Computational calculation
Interesting systems are in general complex, for instance compounds, biomolecules, polymers, etc. Because of the complexity of these systems the analytical calculation of the density of states is in most of the cases impossible. Computer simulations offer a set of algorithms to evaluate the density of states with a high accuracy. One of these algorithms is called the Wang and Landau algorithm.
Within the Wang and Landau scheme any previous knowledge of the density of states is required. One proceeds as follows: the cost function (for example the energy) of the system is discretized. Each time the bin i is reached one updates a histogram for the density of states, , by
where is called the modification factor. As soon as each bin in the histogram is visited a certain number of times (10-15), the modification factor is reduced by some criterion, for instance,
where denotes the -th update step. The simulation finishes when the modification factor is less than a certain threshold, for instance
The Wang and Landau algorithm has some advantages over other common algorithms such as multicanonical simulations and parallel tempering. For example, the density of states is obtained as the main product of the simulation. Additionally, Wang and Landau simulations are completely independent of the temperature. This feature allows to compute the density of states of systems with very rough energy landscape such as proteins.
Mathematically the density of states is formulated in terms of a tower of covering maps.
Local density of states
An important feature of the definition of the DOS is that it can be extended to any system. One of its properties are the translationally invariability which means that the density of the states is homogeneous and it's the same at each point of the system. But this is just a particular case and the LDOS gives a wider description with a heterogeneous density of states through the system.
Concept
Local density of states (LDOS) describes a space-resolved density of states. In materials science, for example, this term is useful when interpreting the data from a scanning tunneling microscope (STM), since this method is capable of imaging electron densities of states with atomic resolution. According to crystal structure, this quantity can be predicted by computational methods, as for example with density functional theory.
A general definition
In a local density of states the contribution of each state is weighted by the density of its wave function at the point. becomes
the factor of means that each state contributes more in the regions where the density is high. An average over of this expression will restore the usual formula for a DOS. The LDOS is useful in inhomogeneous systems, where contains more information than alone.
For a one-dimensional system with a wall, the sine waves give
where .
In a three-dimensional system with the expression is
In fact, we can generalise the local density of states further to
this is called the spectral function and it's a function with each wave function separately in its own variable. In more advanced theory it is connected with the Green's functions and provides a compact representation of some results such as optical absorption.
Solid state devices
LDOS can be used to gain profit into a solid-state device. For example, the figure on the right illustrates LDOS of a transistor as it turns on and off in a ballistic simulation. The LDOS has clear boundary in the source and drain, that corresponds to the location of band edge. In the channel, the DOS is increasing as gate voltage increase and potential barrier goes down.
Optics and photonics
In optics and photonics, the concept of local density of states refers to the states that can be occupied by a photon. For light it is usually measured by fluorescence methods, near-field scanning methods or by cathodoluminescence techniques. Different photonic structures have different LDOS behaviors with different consequences for spontaneous emission. In photonic crystals, near-zero LDOS are expected, inhibiting spontaneous emission.
Similar LDOS enhancement is also expected in plasmonic cavity.
However, in disordered photonic nanostructures, the LDOS behave differently. They fluctuate spatially with their statistics, and are proportional to the scattering strength of the structures.
In addition, the relationship with the mean free path of the scattering is trivial as the LDOS can be still strongly influenced by the short details of strong disorders in the form of a strong Purcell enhancement of the emission. and finally, for the plasmonic disorder, this effect is much stronger for LDOS fluctuations as it can be observed as a strong near-field localization.
| Physical sciences | Statistical mechanics | Physics |
526042 | https://en.wikipedia.org/wiki/Sedna%20%28dwarf%20planet%29 | Sedna (dwarf planet) | Sedna (minor-planet designation: 90377 Sedna) is a dwarf planet in the outermost reaches of the Solar System, orbiting the Sun beyond the orbit of Neptune. Discovered in 2003, the planetoid's surface is one of the reddest known among Solar System bodies. Spectroscopy has revealed Sedna's surface to be mostly a mixture of the solid ices of water, methane, and nitrogen, along with widespread deposits of reddish-colored tholins, a chemical makeup similar to those of some other trans-Neptunian objects. Within the range of uncertainties, it is tied with the dwarf planet in the asteroid belt as the largest dwarf planet not known to have a moon. Its diameter is roughly 1,000 km (most likely in between those of Ceres and Saturn's moon Tethys). Owing to its lack of known moons, the Keplerian laws of planetary motion cannot be employed for determining its mass, and the precise figure remains as yet unknown.
Sedna's orbit is one of the widest known in the Solar System. Its aphelion, the farthest point from the Sun in its elliptical orbit, is located 937 astronomical units (AU) away. This is some 31 times the distance of Neptune's aphelion, and 19 times that of Pluto, spending most of its highly elongated orbit well beyond the heliopause, the boundary beyond which the influence of particles from interstellar space dominates over that of the Sun. Sedna's orbit is also one of the most narrow and elliptical discovered, with an eccentricity of 0.8496. This means that its perihelion, or point of closest approach to the Sun, at 76 AU is around 12.3 times closer than its aphelion. At perihelion, Sedna is only 55% further than Pluto's aphelion. , Sedna is near perihelion, from the Sun, and 2.8 times farther away than Neptune. The dwarf planets and are presently farther away from the Sun than Sedna. It is suggested that an exploratory fly-by mission to Sedna near its perihelion through a Jupiter gravity assist could be completed in 24.5 years.
Due to its exceptionally elongated orbit, the dwarf planet takes approximately 11,400 years, over 11 millennia, to return to the same point in its orbit around the Sun. The International Astronomical Union (IAU) initially considered Sedna to be a member of the scattered disc, a group of objects sent into high-eccentricity orbits by the gravitational influence of Neptune. However, several astronomers who worked in the associated field contested this classification as even its perihelion is far too distant for it to have been scattered by any of the currently known planets. This has led some astronomers to informally refer to it as the first known member of the inner Oort cloud. The dwarf planet is also the prototype of a new orbit class of objects named after itself, the sednoids, which include and Leleākūhonua, all celestial bodies with large perihelion distances and extremely elongated orbits.
The astronomer Michael E. Brown, co-discoverer of Sedna, believes that studying Sedna's unusual orbit could yield valuable information on the origin and early evolution of the Solar System. It might have been perturbed into its orbit by a star within the Sun's birth cluster, or captured from a nearby wandering star, or to have been sent into its present orbit through a close gravitational encounter with the hypothetical 9th planet, sometime during the solar system's formation. The statistically unusual clustering to one side of the solar system of the aphelions of Sedna and other similar objects is speculated to be the evidence for the existence of a planet beyond the orbit of Neptune, which would by itself orbit on the opposing side of the Sun.
History
Discovery
Sedna (provisionally designated ) was discovered by Michael Brown (Caltech), Chad Trujillo (Gemini Observatory), and David Rabinowitz (Yale University) on 14 November 2003. The discovery formed part of a survey begun in 2001 with the Samuel Oschin telescope at Palomar Observatory near San Diego, California, using Yale's 160-megapixel Palomar Quest camera. On that day, an object was observed to move by 4.6 arcseconds over 3.1 hours relative to stars, which indicated that its distance was about 100 AU. Follow-up observations were made in November–December 2003 with the SMARTS (Small and Medium Research Telescope System) at Cerro Tololo Inter-American Observatory in Chile, the Tenagra IV telescope in Nogales, Arizona, and the Keck Observatory on Mauna Kea in Hawaii. Combined with precovery observations taken at the Samuel Oschin telescope in August 2003, and by the Near-Earth Asteroid Tracking consortium in 2001–2002, these observations allowed the accurate determination of its orbit. The calculations showed that the object was moving along a distant and highly eccentric orbit, at a distance of 90.3 AU from the Sun. Precovery images have since been found in the Palomar Digitized Sky Survey dating back to 25 September 1990.
Naming
Brown initially nicknamed Sedna "The Flying Dutchman", or "Dutch", after a legendary ghost ship, because its slow movement had initially masked its presence from his team. He eventually settled on the official name after the goddess Sedna from Inuit mythology, partly because he mistakenly thought the Inuit were the closest polar culture to his home in Pasadena, and partly because the name, unlike Quaoar, would be easily pronounceable by English speakers. Brown further justified his choice of naming by stating that the goddess Sedna's traditional location at the bottom of the Arctic Ocean reflected Sedna's large distance from the Sun. He suggested to the International Astronomical Union's (IAU) Minor Planet Center that any objects discovered in Sedna's orbital region in the future should be named after mythical entities in Arctic mythologies.
The team made the name "Sedna" public before the object had been officially numbered, which caused some controversy among the community of amateur astronomers. Brian Marsden, the head of the Minor Planet Center, stated that such an action was a violation of protocol, and that some members of the IAU might vote against it. Despite the complaints, no objection was raised to the name, and no competing names were suggested. The IAU's Committee on Small Body Nomenclature accepted the name in September 2004, and considered that, in similar cases of extraordinary interest, it might in the future allow names to be announced before they were officially numbered.
Sedna has no symbol in astronomical literature, as the usage of planetary symbols is discouraged in astronomy. Unicode includes a symbol (U+2BF2), but this is mostly used among astrologers. The symbol is a monogram of Sanna, the modern pronunciation of Sedna's name.
Orbit and rotation
Sedna has the longest orbital period of any known object in the Solar System of its size or larger with an orbital period of around 11,400 years. Its orbit is extremely eccentric, with an aphelion of approximately 937 AU and a perihelion of 76.19 AU. Near aphelion, Sedna is one of the coldest places in the Solar System, located far past the termination shock, where temperatures never exceed −240°C (−400°F) due to its extreme distance. At aphelion, Sun as viewed from Sedna is a particularly bright star, among the other stars, in the otherwise black sky, being about 45% as bright as the full moon as seen from Earth. Its perihelion was the largest for any known Solar System object until the discovery of the sednoid . At its aphelion, Sedna orbits the Sun at a meagre 377 m/s, 1.3% that of Earth's average orbital speed.
When Sedna was first discovered, it was 89.6 AU away from the Sun, approaching perihelion, and was the most distant object in the Solar System observed. Sedna was later surpassed by Eris, which was detected by the same survey near its aphelion at 97 AU. Because Sedna is near perihelion , both Eris and are farther from the Sun, at 96 AU and 89 AU respectively, than Sedna at 84 AU, despite both of their semi-major axes being shorter than Sedna's. The orbits of some long-period comets extend further than that of Sedna; they are too dim to be discovered except when approaching perihelion in the inner Solar System. As Sedna nears its perihelion in mid-2076, the Sun will appear merely as a very bright pinpoint in its sky, the G-type star too far away to be visible as a disc to the naked eye.
When first discovered, Sedna was thought to have an unusually long rotational period (20 to 50 days). It was initially speculated that Sedna's rotation was slowed by the gravitational pull of a large binary companion, similar to Pluto's moon Charon. However, a search for such a satellite by the Hubble Space Telescope in March 2004 found no such objects. Subsequent measurements from the MMT telescope showed that Sedna in reality has a much shorter rotation period of about 10 hours, more typical for a body its size. It could rotate in about 18 hours instead, but this is thought to be unlikely.
Physical characteristics
Sedna has a V band absolute magnitude of about 1.8, and is estimated to have an albedo (reflectivity) of around 0.41, giving it a diameter of approximately 900 km. At the time of discovery it was the brightest object found in the Solar System since Pluto in 1930. In 2004, the discoverers placed an upper limit of 1,800 km on its diameter; after observations by the Spitzer Space Telescope, this was revised downward by 2007 to less than 1,600 km. In 2012, measurements from the Herschel Space Observatory suggested that Sedna's diameter was , which would make it smaller than Pluto's moon Charon. In 2013, the same team re-analyzed Sedna's thermal data with an improved thermophysical model and found a consistent value of , suggesting that the original model fit was too precise. Australian observations of a stellar occultation by Sedna in 2013 produced similar results on its diameter, giving chord lengths and . The size of this object suggests it could have undergone differentiation and may have a sub-surface liquid ocean and possibly geologic activity.
As Sedna has no known moons, the direct determination of its mass is as yet impossible without either sending a space probe or perhaps locating a nearby object which is gravitationally perturbed by the planetoid. It is the largest trans-Neptunian Sun-orbiting object not known to have a natural satellite. As of 2024, observations from the Hubble Space Telescope in 2004 have been the only published attempt to find a satellite, and it is possible that a satellite could have been lost in the glare from Sedna itself.
Observations from the SMARTS telescope show that Sedna, in visible light, is one of the reddest objects known in the Solar System, nearly as red as Mars. Its deep red spectral slope is indicative of high concentrations of organic material on its surface. Chad Trujillo and his colleagues suggest that Sedna's dark red color is caused by an extensive surface coating of hydrocarbon sludge, termed tholins. Tholins are a reddish-colored, amorphous, and heterogeneous organic mixture hypothesized to have been transmuted from simpler organic compounds, following billions of years of continuous exposure to ultraviolet radiation, interstellar particles, and other harsh environs as the dwarf planet either comes close to the Sun or transits interstellar space. Its surface is homogeneous in color and spectrum; this may be because Sedna, unlike objects nearer the Sun, is rarely impacted by other bodies, which would expose bright patches of fresh icy material like that on 8405 Asbolus. Sedna and two other very distant objects – and – share their color with outer classical Kuiper belt objects and the centaur 5145 Pholus, suggesting a similar region of origin.
Trujillo and colleagues have placed upper limits on Sedna's surface composition of 60% for methane ice and 70% for water ice. The presence of methane further supports the existence of tholins on Sedna's surface, as methane is among the organic compounds capable of giving rise to tholins. Barucci and colleagues compared Sedna's spectrum with that of Triton and detected weak absorption bands belonging to methane and nitrogen ices. From these observations, they suggested the following model of the surface: 24% Triton-type tholins, 7% amorphous carbon, 10% nitrogen ices, 26% methanol, and 33% methane. The detection of methane and water ice was confirmed in 2006 by the Spitzer Space Telescope mid-infrared photometry. The European Southern Observatory's Very Large Telescope observed Sedna with the SINFONI near-infrared spectrometer, finding indications of tholins and water ice on the surface.
In 2022, low-resolution near-infrared (0.7–5 μm) spectroscopic observations by the James Webb Space Telescope (JWST) revealed the presence of significant amounts of ethane ice (C2H6) and of complex organics on the surface of Sedna. The JWST spectra also contain evidence of the existence of small amounts of ethylene (C2H4), acetylene (C2H2) and possibly carbon dioxide (CO2). On the other hand little evidence of the existence of methane (CH4) and nitrogen ices was found at variance with the earlier observations.
The possible presence of nitrogen on the surface suggests that, at least for a short time, Sedna may have a tenuous atmosphere. During the 200-year portion of its orbit near perihelion, the maximum temperature on Sedna should exceed , the transition temperature between alpha-phase solid N2 and the beta-phase seen on Triton. At 38 K, the N2 vapor pressure would be 14 microbar (1.4 Pa). The weak methane absorption bands indicate that methane on Sedna's surface is ancient, as opposed to being freshly deposited. This finding indicates that Sedna's surface never reaches a temperature high enough for methane on the surface to evaporate and subsequently fall back as snow, which happens on Triton and probably on Pluto.
Origin
In their paper announcing the discovery of Sedna, Brown and his colleagues described it as the first observed body belonging to the Oort cloud, the hypothetical cloud of comet-like objects thought to exist out to nearly a light-year from the Sun. They observed that, unlike scattered disc objects such as Eris, Sedna's perihelion (76 AU) is too distant for it to have been scattered by the gravitational influence of Neptune. Because it is considerably closer to the Sun than was expected for an Oort cloud object, and has an inclination roughly in line with the planets and the Kuiper belt, they described the planetoid as being an "inner Oort cloud object", situated in the disc reaching from the Kuiper belt to the spherical part of the cloud.
If Sedna formed in its current location, the Sun's original protoplanetary disc must have extended as far as 75 AU into space. On top of that, Sedna's initial orbit must have been approximately circular, otherwise its formation by the accretion of smaller bodies into a whole would not have been possible, because the large relative velocities between planetesimals would have been too disruptive. Therefore, it must have been tugged into its current eccentric orbit by a gravitational interaction with another body. In their initial paper, Brown, Rabinowitz and colleagues suggested three possible candidates for the perturbing body: an unseen planet beyond the Kuiper belt, a single passing star, or one of the young stars embedded with the Sun in the stellar cluster in which it formed.
Brown and his team favored the hypothesis that Sedna was lifted into its current orbit by a star from the Sun's birth cluster, arguing that Sedna's aphelion of about 1,000 AU, which is relatively close compared to those of long-period comets, is not distant enough to be affected by passing stars at their current distances from the Sun. They propose that Sedna's orbit is best explained by the Sun having formed in an open cluster of several stars that gradually disassociated over time. That hypothesis has also been advanced by both Alessandro Morbidelli and Scott Jay Kenyon. Computer simulations by Julio A. Fernandez and Adrian Brunini suggest that multiple close passes by young stars in such a cluster would pull many objects into Sedna-like orbits. A study by Morbidelli and Levison suggested that the most likely explanation for Sedna's orbit was that it had been perturbed by a close (approximately 800 AU) pass by another star in the first 100 million years or so of the Solar System's existence.
The trans-Neptunian planet hypothesis has been advanced in several forms by numerous astronomers, including Rodney Gomes and Patryk Lykawka. One scenario involves perturbations of Sedna's orbit by a hypothetical planetary-sized body in the inner Oort cloud. In 2006, simulations suggested that Sedna's orbital traits could be explained by perturbations of a Jupiter-mass () object at 5,000 AU (or less), a Neptune-mass object at 2,000 AU, or even an Earth-mass object at 1,000 AU. Computer simulations by Patryk Lykawka have indicated that Sedna's orbit may have been caused by a body roughly the size of Earth, ejected outward by Neptune early in the Solar System's formation and currently in an elongated orbit between 80 and 170 AU from the Sun. Brown's various sky surveys have not detected any Earth-sized objects out to a distance of about 100 AU. It's a possibility that such an object may have been scattered out of the Solar System after the formation of the inner Oort cloud.
Caltech researchers Konstantin Batygin and Mike Brown have hypothesized the existence of a super-Earth planet in the outer Solar System—Planet Nine—to explain the orbits of a group of extreme trans-Neptunian objects that includes Sedna. This planet would be perhaps six times as massive as Earth. It would have a highly eccentric orbit, and its average distance from the Sun would be about 15 times that of Neptune (which orbits at an average distance of ). Accordingly, its orbital period would be approximately 7,000 to 15,000 years.
Morbidelli and Kenyon have suggested that Sedna did not originate in the Solar System, but was captured by the Sun from a passing extrasolar planetary system, specifically that of a brown dwarf about 1/20th the mass of the Sun () or a main-sequence star 80 percent more massive than the Sun, which, owing to its larger mass, may now be a white dwarf. In either case, the stellar encounter had likely occurred within 100 million years after the Sun's formation. Stellar encounters during this time would have minimal effect on the Oort cloud's final mass and population since the Sun had excess material for replenishing the Oort cloud.
Population
Sedna's highly elliptical orbit, and thus a narrow temporal window for detection and observation with currently available technology, means that the probability of its detection was roughly 1 in 80. Unless its discovery were a fluke, it is expected that another 40–120 Sedna-sized objects with roughly the same orbital parameters would exist in the outer solar system.
In 2007, astronomer Megan Schwamb outlined how each of the proposed mechanisms for Sedna's extreme orbit would affect the structure and dynamics of any wider population. If a trans-Neptunian planet was responsible, all such objects would share roughly the same perihelion (about 80 AU). If Sedna was captured from another planetary system that rotated in the same direction as the Solar System, then all of its population would have orbits on relatively low inclinations and have semi-major axes ranging from 100 to 500 AU. If it rotated in the opposite direction, then two populations would form, one with low and one with high inclinations. The perturbations from passing stars would produce a wide variety of perihelia and inclinations, each dependent on the number and angle of such encounters.
A larger sample of objects with Sedna's extreme perihelion may help in determining which scenario is most likely. "I call Sedna a fossil record of the earliest Solar System", said Brown in 2006. "Eventually, when other fossil records are found, Sedna will help tell us how the Sun formed and the number of stars that were close to the Sun when it formed." A 2007–2008 survey by Brown, Rabinowitz, and Megan Schwamb attempted to locate another member of Sedna's hypothetical population. Although the survey was sensitive to movement out to 1,000 AU and discovered the likely dwarf planet Gonggong, it detected no new sednoid. Subsequent simulations incorporating the new data suggested about 40 Sedna-sized objects probably exist in this region, with the brightest being about Eris's magnitude (−1.0).
In 2014, Chad Trujillo and Scott Sheppard announced the discovery of , an object half the size of Sedna, a 4,200-year orbit similar to Sedna's, and a perihelion within Sedna's range of roughly 80 AU; they speculated that this similarity of orbits may be due to the gravitational shepherding effect of a trans-Neptunian planet. Another high-perihelion trans-Neptunian object was announced by Sheppard and colleagues in 2018, provisionally designated and now named Leleākūhonua. With a perihelion of 65 AU and an even more distant orbit with a period of 40,000 years, its longitude of perihelion (the location where it makes its closest approach to the Sun) appears to be aligned with the directions of both Sedna and , strengthening the case for an apparent orbital clustering of trans-Neptunian objects suspected to be influenced by a hypothetical distant planet, dubbed Planet Nine. In a study detailing Sedna's population and Leleākūhonua's orbital dynamics, Sheppard concluded that the discovery implies a population of about 2 million inner Oort Cloud objects larger than 40 km, with a total mass in the range of (several times the mass of the asteroid belt and 80% the mass of Pluto).
Sedna was recovered from Transiting Exoplanet Survey Satellite data in 2020, as part of preliminary work for an all-sky survey searching for Planet Nine and other as-yet-unknown trans-Neptunian objects.
Classification
The discovery of Sedna renewed the old question of just which astronomical objects ought to be considered planets, and which ones ought not to be. On 15 March 2004, articles on Sedna in the popular press reported misleadingly that a tenth planet had been discovered. This question was resolved for many astronomers by applying the International Astronomical Union's definition of a planet, adopted on 24 August 2006, which mandated that a planet must have cleared the neighborhood around its orbit. Sedna is not expected to have cleared its neighborhood; quantitatively speaking, its Stern–Levison parameter is estimated to be much less than 1. The IAU also adopted dwarf planet as a term for the largest non-planets (despite the name) that, like planets, are in hydrostatic equilibrium and thus can display planet-like geological activity, yet have not cleared their orbital neighborhoods. Sedna is bright enough, and therefore large enough, that it is expected to be in hydrostatic equilibrium. Hence, astronomers generally consider Sedna a dwarf planet.
Besides its physical classification, Sedna is also categorized according to its orbit. The Minor Planet Center, which officially catalogs the objects in the Solar System, designates Sedna only as a trans-Neptunian object (as it orbits beyond Neptune), as does the JPL Small-Body Database. The question of a more precise orbital classification has been much debated, and many astronomers have suggested that the sednoids, together with similar objects such as , be placed in a new category of distant objects named extended scattered disc objects (E-SDO), detached objects, distant detached objects (DDO), or scattered-extended in the formal classification by the Deep Ecliptic Survey.
Exploration
Sedna will come to perihelion around July 2076. This close approach to the Sun provides a window of opportunity for studying it that will not occur again for more than 11 thousand years. Because Sedna spends much of its orbit beyond the heliopause, the point at which the solar wind gives way to the interstellar particle wind, examining Sedna's surface would provide unique information on the effects of interstellar radiation, as well as the properties of the solar wind at its farthest extent. It was calculated in 2011 that a flyby mission to Sedna could take 24.48 years using a Jupiter gravity assist, based on launch dates of 6 May 2033 or 23 June 2046. Sedna would be either 77.27 or 76.43 AU from the Sun when the spacecraft arrives near the end of 2057 or 2070, respectively. Other potential flight trajectories involve gravity assists from Venus, Earth, Saturn, and Neptune as well as Jupiter. Research at the University of Tennessee has also examined the potential for a lander.
| Physical sciences | Solar System | Astronomy |
526164 | https://en.wikipedia.org/wiki/Arabian%20plate | Arabian plate | The Arabian plate is a minor tectonic plate in the Northern and Eastern Hemispheres.
It is one of the three continental plates (along with the African and Indian plates) that have been moving northward in geological history and colliding with the Eurasian plate. This collision is resulting in a mingling of plate pieces and mountain ranges extending in the west from the Pyrenees, crossing Southern Europe to Iranian plateau, forming the Alborz and the Zagros Mountains, to the Himalayas and ranges of Southeast Asia.
Lexicology
The Arabian plate is a designation of the region, and it is also sometimes referred to as the Arab plate.
Borders
The Arabian plate consists mostly of the Arabian Peninsula; it extends westward to the Sinai Peninsula and the Red Sea and northward to the Levant. The plate borders are:
East, with the Indo-Australian plate, at the Owen fracture zone
South, with the African plate to the west and the Somali plate and the Indo-Australian plate to the east
West, a left lateral fault boundary with the African plate called the Dead Sea Transform (DST), and a divergent boundary with the African plate called the Red Sea Rift which runs the length of the Red Sea;
North, convergent boundary with the Anatolian plate and Eurasian plate, including the East Anatolian Fault, Zagros fold and thrust belt, and Makran Trench.
History
The Arabian plate was part of the African plate during most of the Phanerozoic Eon (Paleozoic–Cenozoic), until the Oligocene Epoch of the Cenozoic Era. The Red Sea rifting began in the Eocene, and the separation of Africa and Arabia occurred approximately in the Oligocene, and since then the Arabian plate has been moving toward the Eurasian plate. The opening of the Red Sea rift led to volcanic activity. There are volcanic fields called the Older Harrats, such as Harrat Khaybar and Harrat Rahat, cover parts of the western Arabian plate. Some activity still continues especially around Medina, and there are regular eruptions within the Red Sea.
The collision between the Arabian plate and Eurasia is pushing up the Zagros Mountains of Iran. Because the Arabian plate and Eurasian plate collide, some cities such as those in southeastern Turkey (which is on the Arabian plate) may undergo earthquakes, tsunamis, and volcanoes.
Countries and regions
Countries within the plate include Bahrain, Djibouti, Iraq, Jordan, Kuwait, Oman, Qatar, Saudi Arabia, Syria, United Arab Emirates and Yemen. Regions include the Anti-Lebanon Mountains (Lebanon), parts of Awdal (Somalia/Somaliland), the Khuzestan province (Iran), the Southeastern Anatolia region (Turkey), and the Southern Denkalya subregion (Eritrea).
| Physical sciences | Tectonic plates | Earth science |
20110668 | https://en.wikipedia.org/wiki/Endangered%20species | Endangered species | An endangered species is a species that is very likely to become extinct in the near future, either worldwide or in a particular political jurisdiction. Endangered species may be at risk due to factors such as habitat loss, poaching, invasive species, and climate change. The International Union for Conservation of Nature (IUCN) Red List lists the global conservation status of many species, and various other agencies assess the status of species within particular areas. Many nations have laws that protect conservation-reliant species which, for example, forbid hunting, restrict land development, or create protected areas. Some endangered species are the target of extensive conservation efforts such as captive breeding and habitat restoration.
Human activity is a significant cause in causing some species to become endangered.
Conservation status
The conservation status of a species indicates the likelihood that it will become extinct. Multiple factors are considered when assessing the status of a species; e.g., such statistics as the number remaining, the overall increase or decrease in the population over time, breeding success rates, or known threats. The IUCN Red List of Threatened Species is the best-known worldwide conservation status listing and ranking system.
Over 50% of the world's species are estimated to be at risk of extinction, but the frontier between categories such as 'endangered', 'rare', or 'locally extinct' species is often difficult to draw given the general paucity of data on most of these species. This is notably the case in the world Ocean where endangered species not seen for decades may go extinct unnoticed.
Internationally, 195 countries have signed an accord to create Biodiversity Action Plans that will protect endangered and other threatened species. In the United States, such plans are usually called Species Recovery Plans.
IUCN Red List
Though labeled a list, the IUCN Red List is a system of assessing the global conservation status of species that includes "Data Deficient" (DD) species – species for which more data and assessment is required before their situation may be determined – as well species comprehensively assessed by the IUCN's species assessment process. The species under the index include: mammals, birds, amphibians, cycads, and corals. Those species of "Near Threatened" (NT) and "Least Concern" (LC) status have been assessed and found to have relatively robust and healthy populations, though these may be in decline. Unlike their more general use elsewhere, the List uses the terms "endangered species" and "threatened species" with particular meanings: "Endangered" (EN) species lie between "Vulnerable" (VU) and "Critically Endangered" (CR) species. In 2012, the IUCN Red List listed 3,079 animal and 2,655 plant species as endangered (EN) worldwide.
In Brazil
Brazil is one of the most biodiverse countries in the world, if not the most. It houses not only the Amazon forest but the Atlantic forest, the savanna-like Cerrado among other biomes. Due to the high density of some of its well-preserved rainforests, wildlife trafficking, which along with deforestation is one of the biggest endangerment drivers in Brazil, has become a challenge. Brazil has a broad legal system meant to protect the environment, including its Constitution, as well as several federal, state and local government agencies tasked with protecting the fauna and flora, fining individuals or companies linked to environmental crimes and confiscating illegally taken wildlife. Though such agencies can collect their data, each system operates relatively on its own when it comes to wildlife trafficking. However, both the agencies and the NGO's working in Brazil agree that the birds account for about 80% of trafficked species in the country.
The relation between wildlife smuggling, other environment crimes under the Brazilian law such as deforestation, and endangered species is particularly intricate and troubling since the rarer the animal or plant gets the most targeted and valuable they become in the black market, which leads to more endangered species in its turn.
Additionally, some environment experts and scientists point to the disbanding of environment agencies and the repeal of laws in Brazil under the presidency of Jair Bolsonaro as one of the reasons behind a surge in the number of endangered species. In one occasion during his presidency some fines totaling US$3.1 billion on environment criminals were revoked and at least one fine (related to illegal fishing) imposed on Bolsonaro himself was cancelled and the agent who fined him was demoted.
In the past, Brazil has successfully saved the endemic golden lion tamarin from extinction. Massive campaigns to raise awareness among people by NGO's and governments, which included printing depictions of the golden lion tamarin in the 20 reais Brazilian banknotes (still in circulation), are credited with getting the species out of the critically endangered animals list.
In the United States
There is data from the United States that shows a correlation between human populations and threatened and endangered species. Using species data from the Database on the Economics and Management of Endangered Species database and the period that the Endangered Species Act (ESA) has been in existence, 1970 to 1997, a table was created that suggests a positive relationship between human activity and species endangerment.
Effect of climate change on endangered species
Carbon dioxide in Earth's atmosphere is asserted to be one of the leading causes of animal endangerment. According to the US National Park Service: If we can sufficiently reduce greenhouse gas emissions, many of them will still have a chance to survive and recover. NASA scientist James Hanson has warned that in order to maintain a climate similar to that under which human civilization developed and similar to that which so many organisms are adapted, we need to quickly reduce the carbon dioxide in our atmosphere to 350 parts per million (ppm). Before the industrial revolution, atmospheric carbon dioxide levels rarely rose above 280 ppm; during the 2014 calendar year, carbon dioxide levels fluctuated between 395 and 402 ppm.
Endangered Species Act
Under the Endangered Species Act of 1973 in the United States, species may be listed as "endangered" or "threatened". "The Salt Creek tiger beetle" is an example of an endangered subspecies protected under the ESA. The US Fish and Wildlife Service, as well as the National Marine Fisheries Service are held responsible for classifying and protecting endangered species. They are also responsible for adding a particular species to the list, which can be a long, controversial process.
Some endangered species laws are controversial. Typical areas of controversy include criteria for placing a species on the endangered species list and rules for removing a species from the list once its population has recovered. Whether restrictions on land development constitute a "taking" of land by the government; the related question of whether private landowners should be compensated for the loss of uses of their areas; and obtaining reasonable exceptions to protection laws. Also lobbying from hunters and various industries like the petroleum industry, construction industry, and logging, has been an obstacle in establishing endangered species laws.
The Bush administration lifted a policy that required federal officials to consult a wildlife expert before taking actions that could damage endangered species. Under the Obama administration, this policy was reinstated.
Being listed as an endangered species can have negative effect since it could make a species more desirable for collectors and poachers. This effect is potentially reducible, such as in China where commercially farmed turtles may be reducing some of the pressure to poach endangered species.
Another problem with the listing species is its effect of inciting the use of the "shoot, shovel, and shut-up" method of clearing endangered species from an area of land. Some landowners currently may perceive a diminution in value for their land after finding an endangered animal on it. They have allegedly opted to kill and bury the animals or destroy habitat silently. Thus removing the problem from their land, but at the same time further reducing the population of an endangered species. The effectiveness of the ESA– which coined the term "endangered species"– has been questioned by business advocacy groups and their publications but is nevertheless widely recognized by wildlife scientists who work with the species as an effective recovery tool. Nineteen species have been delisted and recovered and 93% of listed species in the northeastern United States have a recovering or stable population.
Currently, 1,556 endangered species are under protection by government law. This approximation, however, does not take into consideration the species threatened with endangerment that are not included under the protection of laws like the Endangered Species Act. According to NatureServe's global conservation status, approximately thirteen percent of vertebrates (excluding marine fish), seventeen percent of vascular plants, and six to eighteen percent of fungi are considered imperiled. Thus, in total, between seven and eighteen percent of the United States' known animals, fungi and plants are near extinction. This total is substantially more than the number of species protected in the United States under the Endangered Species Act.
Ever since humankind began hunting to preserve itself, over-hunting and fishing have been a large and dangerous problem. Of all the species who became extinct due to interference from humankind, the dodo, passenger pigeon, great auk, Tasmanian tiger and Steller's sea cow are some of the more well known examples; with the bald eagle, grizzly bear, American bison, Eastern timber wolf and sea turtle having been poached to near-extinction. Many began as food sources seen as necessary for survival but became the target of sport. However, due to major efforts to prevent extinction, the bald eagle, or Haliaeetus leucocephalus is now under the category of Least Concern on the red list. A present-day example of the over-hunting of a species can be seen in the oceans as populations of certain whales have been greatly reduced. Large whales like the blue whale, bowhead whale, finback whale, gray whale, sperm whale, and humpback whale are some of the eight whales which are currently still included on the Endangered Species List. Actions have been taken to attempt a reduction in whaling and increase population sizes. The actions include prohibiting all whaling in United States waters, the formation of the CITES treaty which protects all whales, along with the formation of the International Whaling Commission (IWC). But even though all of these movements have been put in place, countries such as Japan continue to hunt and harvest whales under the claim of "scientific purposes". Over-hunting, climatic change and habitat loss leads in landing species in endangered species list. It could mean that extinction rates could increase to a large extent in the future.
In Canada
Endangered species are addressed through Canada's Species at Risk Act. A species is deemed threatened or endangered when it is on the verge of extinction or extirpation. Once a species is deemed threatened or endangered, the Act requires that a recovery plan to be developed that indicates how to stop or reverse the species' population decline. As of 2021, the Committee on the Status of Endangered Wildlife In Canada has assessed 369 species as being endangered in Canada.
In India
The World Wide Fund-India raises concern in the longevity of the following animal species: the Red Panda, the Bengal Tiger, the Ganges River Dolphin, the Asian Elephant.
India signed the Wildlife Protection Act and also joined the Convention on the International Trade in 1976, to prevent poaching from harming its wildlife.
Invasive species
The introduction of non-indigenous species to an area can disrupt the ecosystem to such an extent that native species become endangered. Such introductions may be termed alien or invasive species. In some cases, the invasive species compete with the native species for food or prey on the natives. In other cases, a stable ecological balance may be upset by predation or other causes leading to unexpected species decline. New species may also carry diseases to which the native species have no exposure or resistance.
Climate change
The World Wildlife Fund (WWF) emphasizes that our planet is warming at a rate faster than any time in the past 10,000 years, necessitating species to adapt to new climate patterns, such as variations in rainfall and longer, warmer summers. For example, the U.S. Fish & Wildlife Service highlighted efforts to understand and mitigate the impact of climate change on species through scientific research, modeling, and conservation actions. This includes evaluating the current condition of species, their genetic variation, and how changes in their environment may affect their survival.
The International Union for Conservation of Nature (IUCN) reports that the approximately 1 °C rise in mean global temperature due to human activities is causing serious impacts on species, including changes in abundance, genetic composition, behavior, and survival. The IUCN stresses the importance of environmental policies aimed at reducing CO 2 emissions to lessen the impact of climate change on species. Tools like the IUCN Red List and guidelines for assessing species' vulnerability to climate change are vital for conservation efforts.
In addition, climate change can lead to species decreasing in areas where they once thrived, by being forced to migrate or even going extinct from inhospitable conditions, invasive species, and fragmentation. A study cited by WWF found that one in six species is at risk of extinction due to climate change if no action is taken. The phenomenon of species shifting their ranges in response to changing climates, finding new or shrinking habitats, illustrates the direct impact of global warming on biodiversity. Another major concern is rising ocean acidity caused from excess CO 2 in the atmosphere. This creates acidic conditions in the ocean which creates an inhospitable environment for fish, plants, and other keystone species such as coral reefs
For example the Emperor Penguins, which rely on Antarctic sea ice for breeding, shelter, and food. The melting of ice sheets poses a direct threat to their survival. Similarly, the Mount Rainier white-tailed ptarmigan, adapted to alpine mountaintops, faces habitat loss due to climate changes in snowfall patterns and rising temperatures.
Another example is in the case of the Salton Sea in California. This area is a critical habitat for many endangered and watched species, as well as many migratory birds. Due to environmental shifts from climate change and the addition of agriculture in the surrounding plains, the system has become almost irreparably damaged. The warming temperatures has caused mass evaporation, leaving the Sea much more saline and with much more exposed playa. This not only damages air quality but also has caused fish kills to accumulate as shown pictured below. This has made the system inhospitable to the birds and endangered species relying upon it
Conservation
Captive breeding
Captive breeding is the process of breeding rare or endangered species in human controlled environments with restricted settings, such as wildlife reserves, zoos, and other conservation facilities. Captive breeding is meant to save species from extinction and so stabilise the population of the species that it will not disappear.
This technique has worked for many species for some time, with probably the oldest known such instances of captive mating being attributed to menageries of European and Asian rulers, an example being the Père David's deer. However, captive breeding techniques are usually difficult to implement for such highly mobile species as some migratory birds (e.g. cranes) and fishes (e.g. hilsa). Additionally, if the captive breeding population is too small, then inbreeding may occur due to a reduced gene pool and reduce resistance.In 1981, the Association of Zoos and Aquariums (AZA) created a Species Survival Plan (SSP) to help preserve specific endangered and threatened species through captive breeding. With over 450 SSP Plans, some endangered species are covered by the AZA with plans to cover population management goals and recommendations for breeding for a diverse and healthy population, created by Taxon Advisory Groups. These programs are commonly created as a last resort effort. SSP Programs regularly participate in species recovery, veterinary care for wildlife disease outbreaks, and some other wildlife conservation efforts. The AZA's Species Survival Plan also has breeding and transfer programs, both within and outside of AZA – certified zoos and aquariums. Some animals that are part of SSP programs are giant pandas, lowland gorillas, and California condors.
Private farming
Whereas poaching substantially reduces endangered animal populations, legal, for-profit, private farming does the opposite. It has substantially increased the populations of the southern black rhinoceros and southern white rhinoceros. Richard Emslie, a scientific officer at the IUCN, said of such programs, "Effective law enforcement has become much easier now that the animals are largely privately owned... We have been able to bring local communities into conservation programs. There are increasingly strong economic incentives attached to looking after rhinos rather than simply poaching: from Eco-tourism or selling them on for a profit. So many owners are keeping them secure. The private sector has been key to helping our work."
Conservation experts view the effect of China's turtle farming on the wild turtle populations of China and South-Eastern Asia– many of which are endangered– as "poorly understood". Although they commend the gradual replacement of turtles caught wild with farm-raised turtles in the marketplace– the percentage of farm-raised individuals in the "visible" trade grew from around 30% in 2000 to around 70% in 2007– they worry that many wild animals are caught to provide farmers with breeding stock. The conservation expert Peter Paul van Dijk noted that turtle farmers often believe that animals caught wild are superior breeding stock. Turtle farmers may, therefore, seek and catch the last remaining wild specimens of some endangered turtle species.
In 2015, researchers in Australia managed to coax southern bluefin tuna to breed in landlocked tanks, raising the possibility that fish farming may be able to save the species from overfishing.
Success stories
Hawaiian Monk Seal Rehabilitation: The Hawaiian monk seal are one of the most endangered seal species in the world. Conservation initiatives have focused on mitigating human-seal conflicts, rehabilitating injured seals, and extensive monitoring to ensure their survival. These efforts have led to a gradual increase in their population.
Restoration of the American Bald Eagle: Once on the brink of extinction in the contiguous United States with only 417 known nesting pairs in 1963 due to pesticide use and habitat destruction, the Bald Eagle population has made a remarkable recovery. By 2020, the number of nesting pairs had surged to 71,400. Thanks to habitat protection, legal protection, and DDT ban efforts, leading to the bald eagle being removed from the list of threatened and endangered species.
The Gray Wolf Rebound: Starting in 1995 and 1996, 31 gray wolves from western Canada were relocated to Yellowstone, where they were temporarily kept in acclimation pens before being released into the wild. This careful reintroduction aimed to restore a key predator to the ecosystem, which had profound effects on the park's wildlife dynamics. After being nearly eradicated in the lower 48 states by the early 20th century, reintroduction and protective measures have allowed their populations to rebound significantly. By 2017, gray wolves were delisted in Montana, Idaho, and Wyoming, indicating a recovery to a point where they were no longer considered endangered in these areas.
Recovery of the Channel Island Fox: Beginning in 1999, the Channel Islands National Park launched an ambitious recovery program for the island fox, incorporating several key strategies: captive breeding and reintroduction, removal of predatory golden eagles, re-establishment of bald eagles, and eradication of non-native ungulates. The U.S. Department of the Interior officially recognized the recovery as the fastest for any Endangered Species Act-listed mammal in the U.S., announcing the delisting of three island fox subspecies in 2016. This recovery, from near extinction in the late 1990s to robust populations by the mid-2010s, underscores the power of partnership-driven conservation.
Gallery
| Biology and health sciences | Ecology | null |
20110824 | https://en.wikipedia.org/wiki/Infinity | Infinity | Infinity is something which is boundless, endless, or larger than any natural number. It is denoted by , the infinity symbol.
From the time of the ancient Greeks, the philosophical nature of infinity has been the subject of many discussions among philosophers. In the 17th century, with the introduction of the infinity symbol and the infinitesimal calculus, mathematicians began to work with infinite series and what some mathematicians (including l'Hôpital and Bernoulli) regarded as infinitely small quantities, but infinity continued to be associated with endless processes. As mathematicians struggled with the foundation of calculus, it remained unclear whether infinity could be considered as a number or magnitude and, if so, how this could be done. At the end of the 19th century, Georg Cantor enlarged the mathematical study of infinity by studying infinite sets and infinite numbers, showing that they can be of various sizes. For example, if a line is viewed as the set of all of its points, their infinite number (i.e., the cardinality of the line) is larger than the number of integers. In this usage, infinity is a mathematical concept, and infinite mathematical objects can be studied, manipulated, and used just like any other mathematical object.
The mathematical concept of infinity refines and extends the old philosophical concept, in particular by introducing infinitely many different sizes of infinite sets. Among the axioms of Zermelo–Fraenkel set theory, on which most of modern mathematics can be developed, is the axiom of infinity, which guarantees the existence of infinite sets. The mathematical concept of infinity and the manipulation of infinite sets are widely used in mathematics, even in areas such as combinatorics that may seem to have nothing to do with them. For example, Wiles's proof of Fermat's Last Theorem implicitly relies on the existence of Grothendieck universes, very large infinite sets, for solving a long-standing problem that is stated in terms of elementary arithmetic.
In physics and cosmology, it is an open question whether the universe is spatially infinite or not.
History
Ancient cultures had various ideas about the nature of infinity. The ancient Indians and the Greeks did not define infinity in precise formalism as does modern mathematics, and instead approached infinity as a philosophical concept.
Early Greek
The earliest recorded idea of infinity in Greece may be that of Anaximander (c. 610 – c. 546 BC) a pre-Socratic Greek philosopher. He used the word apeiron, which means "unbounded", "indefinite", and perhaps can be translated as "infinite".
Aristotle (350 BC) distinguished potential infinity from actual infinity, which he regarded as impossible due to the various paradoxes it seemed to produce. It has been argued that, in line with this view, the Hellenistic Greeks had a "horror of the infinite" which would, for example, explain why Euclid (c. 300 BC) did not say that there are an infinity of primes but rather "Prime numbers are more than any assigned multitude of prime numbers." It has also been maintained, that, in proving the infinitude of the prime numbers, Euclid "was the first to overcome the horror of the infinite". There is a similar controversy concerning Euclid's parallel postulate, sometimes translated:
Other translators, however, prefer the translation "the two straight lines, if produced indefinitely ...", thus avoiding the implication that Euclid was comfortable with the notion of infinity. Finally, it has been maintained that a reflection on infinity, far from eliciting a "horror of the infinite", underlay all of early Greek philosophy and that Aristotle's "potential infinity" is an aberration from the general trend of this period.
Zeno: Achilles and the tortoise
Zeno of Elea ( 495 – 430 BC) did not advance any views concerning the infinite. Nevertheless, his paradoxes, especially "Achilles and the Tortoise", were important contributions in that they made clear the inadequacy of popular conceptions. The paradoxes were described by Bertrand Russell as "immeasurably subtle and profound".
Achilles races a tortoise, giving the latter a head start.
Step #1: Achilles runs to the tortoise's starting point while the tortoise walks forward.
Step #2: Achilles advances to where the tortoise was at the end of Step #1 while the tortoise goes yet further.
Step #3: Achilles advances to where the tortoise was at the end of Step #2 while the tortoise goes yet further.
Step #4: Achilles advances to where the tortoise was at the end of Step #3 while the tortoise goes yet further.
Etc.
Apparently, Achilles never overtakes the tortoise, since however many steps he completes, the tortoise remains ahead of him.
Zeno was not attempting to make a point about infinity. As a member of the Eleatics school which regarded motion as an illusion, he saw it as a mistake to suppose that Achilles could run at all. Subsequent thinkers, finding this solution unacceptable, struggled for over two millennia to find other weaknesses in the argument.
Finally, in 1821, Augustin-Louis Cauchy provided both a satisfactory definition of a limit and a proof that, for ,
Suppose that Achilles is running at 10 meters per second, the tortoise is walking at 0.1 meters per second, and the latter has a 100-meter head start. The duration of the chase fits Cauchy's pattern with and . Achilles does overtake the tortoise; it takes him
Early Indian
The Jain mathematical text Surya Prajnapti (c. 4th–3rd century BCE) classifies all numbers into three sets: enumerable, innumerable, and infinite. Each of these was further subdivided into three orders:
Enumerable: lowest, intermediate, and highest
Innumerable: nearly innumerable, truly innumerable, and innumerably innumerable
Infinite: nearly infinite, truly infinite, infinitely infinite
17th century
In the 17th century, European mathematicians started using infinite numbers and infinite expressions in a systematic fashion. In 1655, John Wallis first used the notation for such a number in his De sectionibus conicis, and exploited it in area calculations by dividing the region into infinitesimal strips of width on the order of But in Arithmetica infinitorum (1656), he indicates infinite series, infinite products and infinite continued fractions by writing down a few terms or factors and then appending "&c.", as in "1, 6, 12, 18, 24, &c."
In 1699, Isaac Newton wrote about equations with an infinite number of terms in his work De analysi per aequationes numero terminorum infinitas.
Mathematics
Hermann Weyl opened a mathematico-philosophic address given in 1930 with:
Symbol
The infinity symbol (sometimes called the lemniscate) is a mathematical symbol representing the concept of infinity. The symbol is encoded in Unicode at and in LaTeX as \infty.
It was introduced in 1655 by John Wallis, and since its introduction, it has also been used outside mathematics in modern mysticism and literary symbology.
Calculus
Gottfried Leibniz, one of the co-inventors of infinitesimal calculus, speculated widely about infinite numbers and their use in mathematics. To Leibniz, both infinitesimals and infinite quantities were ideal entities, not of the same nature as appreciable quantities, but enjoying the same properties in accordance with the Law of continuity.
Real analysis
In real analysis, the symbol , called "infinity", is used to denote an unbounded limit. The notation means that increases without bound, and means that decreases without bound. For example, if for every , then
means that does not bound a finite area from to
means that the area under is infinite.
means that the total area under is finite, and is equal to
Infinity can also be used to describe infinite series, as follows:
means that the sum of the infinite series converges to some real value
means that the sum of the infinite series properly diverges to infinity, in the sense that the partial sums increase without bound.
In addition to defining a limit, infinity can be also used as a value in the extended real number system. Points labeled and can be added to the topological space of the real numbers, producing the two-point compactification of the real numbers. Adding algebraic properties to this gives us the extended real numbers. We can also treat and as the same, leading to the one-point compactification of the real numbers, which is the real projective line. Projective geometry also refers to a line at infinity in plane geometry, a plane at infinity in three-dimensional space, and a hyperplane at infinity for general dimensions, each consisting of points at infinity.
Complex analysis
In complex analysis the symbol , called "infinity", denotes an unsigned infinite limit. The expression means that the magnitude of grows beyond any assigned value. A point labeled can be added to the complex plane as a topological space giving the one-point compactification of the complex plane. When this is done, the resulting space is a one-dimensional complex manifold, or Riemann surface, called the extended complex plane or the Riemann sphere. Arithmetic operations similar to those given above for the extended real numbers can also be defined, though there is no distinction in the signs (which leads to the one exception that infinity cannot be added to itself). On the other hand, this kind of infinity enables division by zero, namely for any nonzero complex number . In this context, it is often useful to consider meromorphic functions as maps into the Riemann sphere taking the value of at the poles. The domain of a complex-valued function may be extended to include the point at infinity as well. One important example of such functions is the group of Möbius transformations (see Möbius transformation § Overview).
Nonstandard analysis
The original formulation of infinitesimal calculus by Isaac Newton and Gottfried Leibniz used infinitesimal quantities. In the second half of the 20th century, it was shown that this treatment could be put on a rigorous footing through various logical systems, including smooth infinitesimal analysis and nonstandard analysis. In the latter, infinitesimals are invertible, and their inverses are infinite numbers. The infinities in this sense are part of a hyperreal field; there is no equivalence between them as with the Cantorian transfinites. For example, if H is an infinite number in this sense, then H + H = 2H and H + 1 are distinct infinite numbers. This approach to non-standard calculus is fully developed in .
Set theory
A different form of "infinity" is the ordinal and cardinal infinities of set theory—a system of transfinite numbers first developed by Georg Cantor. In this system, the first transfinite cardinal is aleph-null (ℵ0), the cardinality of the set of natural numbers. This modern mathematical conception of the quantitative infinite developed in the late 19th century from works by Cantor, Gottlob Frege, Richard Dedekind and others—using the idea of collections or sets.
Dedekind's approach was essentially to adopt the idea of one-to-one correspondence as a standard for comparing the size of sets, and to reject the view of Galileo (derived from Euclid) that the whole cannot be the same size as the part. (However, see Galileo's paradox where Galileo concludes that positive integers cannot be compared to the subset of positive square integers since both are infinite sets.) An infinite set can simply be defined as one having the same size as at least one of its proper parts; this notion of infinity is called Dedekind infinite. The diagram to the right gives an example: viewing lines as infinite sets of points, the left half of the lower blue line can be mapped in a one-to-one manner (green correspondences) to the higher blue line, and, in turn, to the whole lower blue line (red correspondences); therefore the whole lower blue line and its left half have the same cardinality, i.e. "size".
Cantor defined two kinds of infinite numbers: ordinal numbers and cardinal numbers. Ordinal numbers characterize well-ordered sets, or counting carried on to any stopping point, including points after an infinite number have already been counted. Generalizing finite and (ordinary) infinite sequences which are maps from the positive integers leads to mappings from ordinal numbers to transfinite sequences. Cardinal numbers define the size of sets, meaning how many members they contain, and can be standardized by choosing the first ordinal number of a certain size to represent the cardinal number of that size. The smallest ordinal infinity is that of the positive integers, and any set which has the cardinality of the integers is countably infinite. If a set is too large to be put in one-to-one correspondence with the positive integers, it is called uncountable. Cantor's views prevailed and modern mathematics accepts actual infinity as part of a consistent and coherent theory. Certain extended number systems, such as the hyperreal numbers, incorporate the ordinary (finite) numbers and infinite numbers of different sizes.
Cardinality of the continuum
One of Cantor's most important results was that the cardinality of the continuum is greater than that of the natural numbers ; that is, there are more real numbers than natural numbers . Namely, Cantor showed that .
The continuum hypothesis states that there is no cardinal number between the cardinality of the reals and the cardinality of the natural numbers, that is, .This hypothesis cannot be proved or disproved within the widely accepted Zermelo–Fraenkel set theory, even assuming the Axiom of Choice.
Cardinal arithmetic can be used to show not only that the number of points in a real number line is equal to the number of points in any segment of that line, but also that this is equal to the number of points on a plane and, indeed, in any finite-dimensional space.
The first of these results is apparent by considering, for instance, the tangent function, which provides a one-to-one correspondence between the interval () and.The second result was proved by Cantor in 1878, but only became intuitively apparent in 1890, when Giuseppe Peano introduced the space-filling curves, curved lines that twist and turn enough to fill the whole of any square, or cube, or hypercube, or finite-dimensional space. These curves can be used to define a one-to-one correspondence between the points on one side of a square and the points in the square.
Geometry
Until the end of the 19th century, infinity was rarely discussed in geometry, except in the context of processes that could be continued without any limit. For example, a line was what is now called a line segment, with the proviso that one can extend it as far as one wants; but extending it infinitely was out of the question. Similarly, a line was usually not considered to be composed of infinitely many points but was a location where a point may be placed. Even if there are infinitely many possible positions, only a finite number of points could be placed on a line. A witness of this is the expression "the locus of a point that satisfies some property" (singular), where modern mathematicians would generally say "the set of the points that have the property" (plural).
One of the rare exceptions of a mathematical concept involving actual infinity was projective geometry, where points at infinity are added to the Euclidean space for modeling the perspective effect that shows parallel lines intersecting "at infinity". Mathematically, points at infinity have the advantage of allowing one to not consider some special cases. For example, in a projective plane, two distinct lines intersect in exactly one point, whereas without points at infinity, there are no intersection points for parallel lines. So, parallel and non-parallel lines must be studied separately in classical geometry, while they need not be distinguished in projective geometry.
Before the use of set theory for the foundation of mathematics, points and lines were viewed as distinct entities, and a point could be located on a line. With the universal use of set theory in mathematics, the point of view has dramatically changed: a line is now considered as the set of its points, and one says that a point belongs to a line instead of is located on a line (however, the latter phrase is still used).
In particular, in modern mathematics, lines are infinite sets.
Infinite dimension
The vector spaces that occur in classical geometry have always a finite dimension, generally two or three. However, this is not implied by the abstract definition of a vector space, and vector spaces of infinite dimension can be considered. This is typically the case in functional analysis where function spaces are generally vector spaces of infinite dimension.
In topology, some constructions can generate topological spaces of infinite dimension. In particular, this is the case of iterated loop spaces.
Fractals
The structure of a fractal object is reiterated in its magnifications. Fractals can be magnified indefinitely without losing their structure and becoming "smooth"; they have infinite perimeters and can have infinite or finite areas. One such fractal curve with an infinite perimeter and finite area is the Koch snowflake.
Mathematics without infinity
Leopold Kronecker was skeptical of the notion of infinity and how his fellow mathematicians were using it in the 1870s and 1880s. This skepticism was developed in the philosophy of mathematics called finitism, an extreme form of mathematical philosophy in the general philosophical and mathematical schools of constructivism and intuitionism.
Physics
In physics, approximations of real numbers are used for continuous measurements and natural numbers are used for discrete measurements (i.e., counting). Concepts of infinite things such as an infinite plane wave exist, but there are no experimental means to generate them.
Cosmology
The first published proposal that the universe is infinite came from Thomas Digges in 1576. Eight years later, in 1584, the Italian philosopher and astronomer Giordano Bruno proposed an unbounded universe in On the Infinite Universe and Worlds: "Innumerable suns exist; innumerable earths revolve around these suns in a manner similar to the way the seven planets revolve around our sun. Living beings inhabit these worlds."
Cosmologists have long sought to discover whether infinity exists in our physical universe: Are there an infinite number of stars? Does the universe have infinite volume? Does space "go on forever"? This is still an open question of cosmology. The question of being infinite is logically separate from the question of having boundaries. The two-dimensional surface of the Earth, for example, is finite, yet has no edge. By travelling in a straight line with respect to the Earth's curvature, one will eventually return to the exact spot one started from. The universe, at least in principle, might have a similar topology. If so, one might eventually return to one's starting point after travelling in a straight line through the universe for long enough.
The curvature of the universe can be measured through multipole moments in the spectrum of the cosmic background radiation. To date, analysis of the radiation patterns recorded by the WMAP spacecraft hints that the universe has a flat topology. This would be consistent with an infinite physical universe.
However, the universe could be finite, even if its curvature is flat. An easy way to understand this is to consider two-dimensional examples, such as video games where items that leave one edge of the screen reappear on the other. The topology of such games is toroidal and the geometry is flat. Many possible bounded, flat possibilities also exist for three-dimensional space.
The concept of infinity also extends to the multiverse hypothesis, which, when explained by astrophysicists such as Michio Kaku, posits that there are an infinite number and variety of universes. Also, cyclic models posit an infinite amount of Big Bangs, resulting in an infinite variety of universes after each Big Bang event in an infinite cycle.
Logic
In logic, an infinite regress argument is "a distinctively philosophical kind of argument purporting to show that a thesis is defective because it generates an infinite series when either (form A) no such series exists or (form B) were it to exist, the thesis would lack the role (e.g., of justification) that it is supposed to play."
Computing
The IEEE floating-point standard (IEEE 754) specifies a positive and a negative infinity value (and also indefinite values). These are defined as the result of arithmetic overflow, division by zero, and other exceptional operations.
Some programming languages, such as Java and J, allow the programmer an explicit access to the positive and negative infinity values as language constants. These can be used as greatest and least elements, as they compare (respectively) greater than or less than all other values. They have uses as sentinel values in algorithms involving sorting, searching, or windowing.
In languages that do not have greatest and least elements but do allow overloading of relational operators, it is possible for a programmer to create the greatest and least elements. In languages that do not provide explicit access to such values from the initial state of the program but do implement the floating-point data type, the infinity values may still be accessible and usable as the result of certain operations.
In programming, an infinite loop is a loop whose exit condition is never satisfied, thus executing indefinitely.
Arts, games, and cognitive sciences
Perspective artwork uses the concept of vanishing points, roughly corresponding to mathematical points at infinity, located at an infinite distance from the observer. This allows artists to create paintings that realistically render space, distances, and forms. Artist M.C. Escher is specifically known for employing the concept of infinity in his work in this and other ways.
Variations of chess played on an unbounded board are called infinite chess.
Cognitive scientist George Lakoff considers the concept of infinity in mathematics and the sciences as a metaphor. This perspective is based on the basic metaphor of infinity (BMI), defined as the ever-increasing sequence <1, 2, 3, …>.
| Mathematics | Analysis | null |
306364 | https://en.wikipedia.org/wiki/Liquid%20nitrogen | Liquid nitrogen | Liquid nitrogen (LN2) is nitrogen in a liquid state at low temperature. Liquid nitrogen has a boiling point of about . It is produced industrially by fractional distillation of liquid air. It is a colorless, mobile liquid whose viscosity is about one-tenth that of acetone (i.e. roughly one-thirtieth that of water at room temperature). Liquid nitrogen is widely used as a coolant.
Physical properties
The diatomic character of the N2 molecule is retained after liquefaction. The weak van der Waals interaction between the N2 molecules results in little interatomic attraction. This is the cause of nitrogen's unusually low boiling point.
The temperature of liquid nitrogen can readily be reduced to its freezing point by placing it in a vacuum chamber pumped by a vacuum pump. Liquid nitrogen's efficiency as a coolant is limited by the fact that it boils immediately on contact with a warmer object, enveloping the object in an insulating layer of nitrogen gas bubbles. This effect, known as the Leidenfrost effect, occurs when any liquid comes in contact with a surface which is significantly hotter than its boiling point. Faster cooling may be obtained by plunging an object into a slush of liquid and solid nitrogen rather than liquid nitrogen alone.
Handling
As a cryogenic fluid that rapidly freezes living tissue, its handling and storage require thermal insulation. It can be stored and transported in vacuum flasks, the temperature being held constant at 77 K by slow boiling of the liquid. Depending on the size and design, the holding time of vacuum flasks ranges from a few hours to a few weeks. The development of pressurised super-insulated vacuum vessels has enabled liquid nitrogen to be stored and transported over longer time periods with losses reduced to 2 percent per day or less.
Uses
Liquid nitrogen is a compact and readily transported source of dry nitrogen gas, as it does not require pressurization. Further, its ability to maintain temperatures far below the freezing point of water, specific heat of 1040 J⋅kg−1⋅K−1 and heat of vaporization of 200 kJ⋅kg−1 makes it extremely useful in a wide range of applications, primarily as an open-cycle refrigerant, including:
in cryotherapy for removing unsightly or potentially malignant skin lesions such as warts and actinic keratosis
to store cells at low temperature for laboratory work
in cryogenics
as a backup nitrogen source in hypoxic air fire prevention systems
as a source of very dry nitrogen gas
for the immersion, freezing, and transportation of food products
for the cryopreservation of blood, reproductive cells (sperm and egg), and other biological samples and materials
to preserve tissue samples from surgical excisions for future studies
to facilitate cryoconservation of animal genetic resources
to freeze water and oil pipes in order to work on them in situations where a valve is not available to block fluid flow to the work area; this method is known as a cryogenic isolation
for cryonic preservation in hopes of future reanimation
for shrink-fitting machinery parts together
as a coolant
for charge-coupled cameras in astronomy
for a high-temperature superconductor to a temperature sufficient to achieve superconductivity
to maintain a low temperature around the primary liquid helium cooling system of high-field superconducting magnets used in e.g. nuclear magnetic resonance spectrometers and magnetic resonance imaging systems
for vacuum pump traps and in controlled-evaporation processes in chemistry
as a component of cooling baths used for very low temperature reactions in chemistry
to increase the sensitivity of infrared homing seeker heads of missiles such as the Strela 3
to temporarily shrink mechanical components during machine assembly and allow improved interference fits
for computers and extreme overclocking
for simulation of space background in vacuum chamber during spacecraft thermal testing
to control the temperature of mass concrete, injected to precool concrete mixes during delivery
in food preparation, such as for making ultra-smooth ice cream. | Physical sciences | Group 15 | Chemistry |
306543 | https://en.wikipedia.org/wiki/Structural%20analysis | Structural analysis | Structural analysis is a branch of solid mechanics which uses simplified models for solids like bars, beams and shells for engineering decision making. Its main objective is to determine the effect of loads on physical structures and their components. In contrast to theory of elasticity, the models used in structural analysis are often differential equations in one spatial variable. Structures subject to this type of analysis include all that must withstand loads, such as buildings, bridges, aircraft and ships. Structural analysis uses ideas from applied mechanics, materials science and applied mathematics to compute a structure's deformations, internal forces, stresses, support reactions, velocity, accelerations, and stability. The results of the analysis are used to verify a structure's fitness for use, often precluding physical tests. Structural analysis is thus a key part of the engineering design of structures.
Structures and loads
In the context to structural analysis, a structure refers to a body or system of connected parts used to support a load. Important examples related to Civil Engineering include buildings, bridges, and towers; and in other branches of engineering, ship and aircraft frames, tanks, pressure vessels, mechanical systems, and electrical supporting structures are important. To design a structure, an engineer must account for its safety, aesthetics, and serviceability, while considering economic and environmental constraints. Other branches of engineering work on a wide variety of non-building structures.
Classification of structures
A structural system is the combination of structural elements and their materials. It is important for a structural engineer to be able to classify a structure by either its form or its function, by recognizing the various elements composing that structure.
The structural elements guiding the systemic forces through the materials are not only such as a connecting rod, a truss, a beam, or a column, but also a cable, an arch, a cavity or channel, and even an angle, a surface structure, or a frame.
Loads
Once the dimensional requirement for a structure have been defined, it becomes necessary to determine the loads the structure must support. Structural design, therefore begins with specifying loads that act on the structure. The design loading for a structure is often specified in building codes. There are two types of codes: general building codes and design codes, engineers must satisfy all of the code's requirements in order for the structure to remain reliable.
There are two types of loads that structure engineering must encounter in the design. The first type of loads are dead loads that consist of the weights of the various structural members and the weights of any objects that are permanently attached to the structure. For example, columns, beams, girders, the floor slab, roofing, walls, windows, plumbing, electrical fixtures, and other miscellaneous attachments. The second type of loads are live loads which vary in their magnitude and location. There are many different types of live loads like building loads, highway bridge loads, railroad bridge loads, impact loads, wind loads, snow loads, earthquake loads, and other natural loads.
Analytical methods
To perform an accurate analysis a structural engineer must determine information such as structural loads, geometry, support conditions, and material properties. The results of such an analysis typically include support reactions, stresses and displacements. This information is then compared to criteria that indicate the conditions of failure. Advanced structural analysis may examine dynamic response, stability and non-linear behavior.
There are three approaches to the analysis: the mechanics of materials approach (also known as strength of materials), the elasticity theory approach (which is actually a special case of the more general field of continuum mechanics), and the finite element approach. The first two make use of analytical formulations which apply mostly simple linear elastic models, leading to closed-form solutions, and can often be solved by hand. The finite element approach is actually a numerical method for solving differential equations generated by theories of mechanics such as elasticity theory and strength of materials. However, the finite-element method depends heavily on the processing power of computers and is more applicable to structures of arbitrary size and complexity.
Regardless of approach, the formulation is based on the same three fundamental relations: equilibrium, constitutive, and compatibility. The solutions are approximate when any of these relations are only approximately satisfied, or only an approximation of reality.
Limitations
Each method has noteworthy limitations. The method of mechanics of materials is limited to very simple structural elements under relatively simple loading conditions. The structural elements and loading conditions allowed, however, are sufficient to solve many useful engineering problems. The theory of elasticity allows the solution of structural elements of general geometry under general loading conditions, in principle. Analytical solution, however, is limited to relatively simple cases. The solution of elasticity problems also requires the solution of a system of partial differential equations, which is considerably more mathematically demanding than the solution of mechanics of materials problems, which require at most the solution of an ordinary differential equation. The finite element method is perhaps the most restrictive and most useful at the same time. This method itself relies upon other structural theories (such as the other two discussed here) for equations to solve. It does, however, make it generally possible to solve these equations, even with highly complex geometry and loading conditions, with the restriction that there is always some numerical error. Effective and reliable use of this method requires a solid understanding of its limitations.
Strength of materials methods (classical methods)
The simplest of the three methods here discussed, the mechanics of materials method is available for simple structural members subject to specific loadings such as axially loaded bars, prismatic beams in a state of pure bending, and circular shafts subject to torsion. The solutions can under certain conditions be superimposed using the superposition principle to analyze a member undergoing combined loading. Solutions for special cases exist for common structures such as thin-walled pressure vessels.
For the analysis of entire systems, this approach can be used in conjunction with statics, giving rise to the method of sections and method of joints for truss analysis, moment distribution method for small rigid frames, and portal frame and cantilever method for large rigid frames. Except for moment distribution, which came into use in the 1930s, these methods were developed in their current forms in the second half of the nineteenth century. They are still used for small structures and for preliminary design of large structures.
The solutions are based on linear isotropic infinitesimal elasticity and Euler–Bernoulli beam theory. In other words, they contain the assumptions (among others) that the materials in question are elastic, that stress is related linearly to strain, that the material (but not the structure) behaves identically regardless of direction of the applied load, that all deformations are small, and that beams are long relative to their depth. As with any simplifying assumption in engineering, the more the model strays from reality, the less useful (and more dangerous) the result.
Example
There are 2 commonly used methods to find the truss element forces, namely the method of joints and the method of sections. Below is an example that is solved using both of these methods. The first diagram below is the presented problem for which the truss element forces have to be found. The second diagram is the loading diagram and contains the reaction forces from the joints.
Since there is a pin joint at A, it will have 2 reaction forces. One in the x direction and the other in the y direction. At point B, there is a roller joint and hence only 1 reaction force in the y direction. Assuming these forces to be in their respective positive directions (if they are not in the positive directions, the value will be negative).
Since the system is in static equilibrium, the sum of forces in any direction is zero and the sum of moments about any point is zero.
Therefore, the magnitude and direction of the reaction forces can be calculated.
Method of joints
This type of method uses the force balance in the x and y directions at each of the joints in the truss structure.
At A,
At D,
At C,
Although the forces in each of the truss elements are found, it is a good practice to verify the results by completing the remaining force balances.
At B,
Method of sections
This method can be used when the truss element forces of only a few members are to be found. This method is used by introducing a single straight line cutting through the member whose force has to be calculated. However this method has a limit in that the cutting line can pass through a maximum of only 3 members of the truss structure. This restriction is because this method uses the force balances in the x and y direction and the moment balance, which gives a maximum of 3 equations to find a maximum of 3 unknown truss element forces through which this cut is made. Find the forces FAB, FBD and FCD in the above example
Method 1: Ignore the right side
Method 2: Ignore the left side
The truss elements forces in the remaining members can be found by using the above method with a section passing through the remaining members.
Elasticity methods
Elasticity methods are available generally for an elastic solid of any shape. Individual members such as beams, columns, shafts, plates and shells may be modeled. The solutions are derived from the equations of linear elasticity. The equations of elasticity are a system of 15 partial differential equations. Due to the nature of the mathematics involved, analytical solutions may only be produced for relatively simple geometries. For complex geometries, a numerical solution method such as the finite element method is necessary.
Methods using numerical approximation
It is common practice to use approximate solutions of differential equations as the basis for structural analysis. This is usually done using numerical approximation techniques. The most commonly used numerical approximation in structural analysis is the Finite Element Method.
The finite element method approximates a structure as an assembly of elements or components with various forms of connection between them and each element of which has an associated stiffness. Thus, a continuous system such as a plate or shell is modeled as a discrete system with a finite number of elements interconnected at finite number of nodes and the overall stiffness is the result of the addition of the stiffness of the various elements. The behaviour of individual elements is characterized by the element's stiffness (or flexibility) relation. The assemblage of the various stiffness's into a master stiffness matrix that represents the entire structure leads to the system's stiffness or flexibility relation. To establish the stiffness (or flexibility) of a particular element, we can use the mechanics of materials approach for simple one-dimensional bar elements, and the elasticity approach for more complex two- and three-dimensional elements. The analytical and computational development are best effected throughout by means of matrix algebra, solving partial differential equations.
Early applications of matrix methods were applied to articulated frameworks with truss, beam and column elements; later and more advanced matrix methods, referred to as "finite element analysis", model an entire structure with one-, two-, and three-dimensional elements and can be used for articulated systems together with continuous systems such as a pressure vessel, plates, shells, and three-dimensional solids. Commercial computer software for structural analysis typically uses matrix finite-element analysis, which can be further classified into two main approaches: the displacement or stiffness method and the force or flexibility method. The stiffness method is the most popular by far thanks to its ease of implementation as well as of formulation for advanced applications. The finite-element technology is now sophisticated enough to handle just about any system as long as sufficient computing power is available. Its applicability includes, but is not limited to, linear and non-linear analysis, solid and fluid interactions, materials that are isotropic, orthotropic, or anisotropic, and external effects that are static, dynamic, and environmental factors. This, however, does not imply that the computed solution will automatically be reliable because much depends on the model and the reliability of the data input.
Timeline
1452–1519 Leonardo da Vinci made many contributions
1638: Galileo Galilei published the book "Two New Sciences" in which he examined the failure of simple structures
1660: Hooke's law by Robert Hooke
1687: Isaac Newton published "Philosophiae Naturalis Principia Mathematica" which contains the Newton's laws of motion
1750: Euler–Bernoulli beam equation
1700–1782: Daniel Bernoulli introduced the principle of virtual work
1707–1783: Leonhard Euler developed the theory of buckling of columns
1826: Claude-Louis Navier published a treatise on the elastic behaviors of structures
1873: Carlo Alberto Castigliano presented his dissertation "Intorno ai sistemi elastici", which contains his theorem for computing displacement as partial derivative of the strain energy. This theorem includes the method of 'least work' as a special case
1878-1972 Stephen Timoshenko father of modern Applied mechanics including the Timoshenko–Ehrenfest beam theory
1936: Hardy Cross' publication of the moment distribution method which was later recognized as a form of the relaxation method applicable to the problem of flow in pipe-network
1941: Alexander Hrennikoff submitted his D.Sc. thesis in MIT on the discretization of plane elasticity problems using a lattice framework
1942: R. Courant divided a domain into finite subregions
1956: J. Turner, R. W. Clough, H. C. Martin, and L. J. Topp's paper on the "Stiffness and Deflection of Complex Structures" introduces the name "finite-element method" and is widely recognized as the first comprehensive treatment of the method as it is known today
| Technology | Disciplines | null |
306609 | https://en.wikipedia.org/wiki/Group%203%20element | Group 3 element | |-
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Group 3 is the first group of transition metals in the periodic table. This group is closely related to the rare-earth elements. It contains the four elements scandium (Sc), yttrium (Y), lutetium (Lu), and lawrencium (Lr). The group is also called the scandium group or scandium family after its lightest member.
The chemistry of the group 3 elements is typical for early transition metals: they all essentially have only the group oxidation state of +3 as a major one, and like the preceding main-group metals are quite electropositive and have a less rich coordination chemistry. Due to the effects of the lanthanide contraction, yttrium and lutetium are very similar in properties. Yttrium and lutetium have essentially the chemistry of the heavy lanthanides, but scandium shows several differences due to its small size. This is a similar pattern to those of the early transition metal groups, where the lightest element is distinct from the very similar next two.
All the group 3 elements are rather soft, silvery-white metals, although their hardness increases with atomic number. They quickly tarnish in air and react with water, though their reactivity is masked by the formation of an oxide layer. The first three of them occur naturally, and especially yttrium and lutetium are almost invariably associated with the lanthanides due to their similar chemistry. Lawrencium is strongly radioactive: it does not occur naturally and must be produced by artificial synthesis, but its observed and theoretically predicted properties are consistent with it being a heavier homologue of lutetium. None of the group 3 elements have any biological role.
Historically, sometimes lanthanum (La) and actinium (Ac) were included in the group instead of lutetium and lawrencium, because the electron configurations of many of the rare earths were initially measured wrongly. This version of group 3 is still commonly found in textbooks, but most authors focusing on the subject are against it. Some authors attempt to compromise between the two formats by leaving the spaces below yttrium blank, but this contradicts quantum mechanics as it results in an f-block that is 15 elements wide rather than 14 (the maximum occupancy of an f-subshell).
Composition
Physical, chemical, and electronic evidence overwhelmingly shows that the correct elements in group 3 are scandium, yttrium, lutetium, and lawrencium: this is the classification adopted by most chemists and physicists who have considered the matter. It was supported by IUPAC in a 1988 report and reaffirmed in 2021. Many textbooks however show group 3 as containing scandium, yttrium, lanthanum, and actinium, a format based on historically wrongly measured electron configurations: Lev Landau and Evgeny Lifshitz already considered it to be "incorrect" in 1948, but the issue was brought to a wide debate only in 1982 by William B. Jensen.
The spaces below yttrium are sometimes left blank as a third option, but there is confusion in the literature on whether this format implies that group 3 contains only scandium and yttrium, or if it also contains all the lanthanides and actinides; either way, this format contradicts quantum physics by creating a 15-element-wide f-block when only 14 electrons can fit in an f-subshell. While the 2021 IUPAC report noted that 15-element-wide f-blocks are supported by some practitioners of a specialised branch of relativistic quantum mechanics focusing on the properties of superheavy elements, the project's opinion was that such interest-dependent concerns should not have any bearing on how the periodic table is presented to "the general chemical and scientific community". In fact, relativistic quantum-mechanical calculations of Lu and Lr compounds found no valence f-orbitals in either element. Other authors focusing on superheavy elements since clarified that the "15th entry of the f-block represents the first slot of the d-block which is left vacant to indicate the place of the f-block inserts", which would imply that this form still has Lu and Lr (the 15th entries in question) as d-block elements under Sc and Y. Indeed, when IUPAC publications expand the table to 32 columns, they make this clear and place Lu and Lr under Y.
As noted by the 2021 IUPAC report, Sc-Y-Lu-Lr is the only form that simultaneously allows for the preservation of the sequence of atomic number, avoids splitting the d-block into "two highly uneven portions", and gives the blocks the correct widths quantum mechanics demands (2, 6, 10, and 14). While arguments in favour of Sc-Y-La-Ac can still be found in the literature, many authors consider them to be logically inconsistent. For example, it has been argued that lanthanum and actinium cannot be f-block elements because their atoms have not begun to fill the f-subshells. But the same is true of thorium which is never disputed as an f-block element, and this argument overlooks the problem on the other end: that the f-shells complete filling at ytterbium and nobelium (matching the Sc-Y-Lu-Lr form), not at lutetium and lawrencium (as in Sc-Y-La-Ac). Lanthanum, actinium, and thorium are simply examples of exceptions to the Madelung rule; not only do those exceptions represent a minority of elements (only 20 out of 118), but they have also never been considered as relevant for positioning any other elements on the periodic table. In gaseous atoms, the d-shells complete their filling at copper (3d104s1), palladium (4d105s0), and gold (5d106s1), but it is universally accepted by chemists that these configurations are exceptional and that the d-block really ends in accordance with the Madelung rule at zinc (3d104s2), cadmium (4d105s2), and mercury (5d106s2). The relevant fact for placement is that lanthanum and actinium (like thorium) have valence f-orbitals that can become occupied in chemical environments, whereas lutetium and lawrencium do not: their f-shells are in the core, and cannot be used for chemical reactions. Thus the relationship between yttrium and lanthanum is only a secondary relationship between elements with the same number of valence electrons but different kinds of valence orbitals, such as that between chromium and uranium; whereas the relationship between yttrium and lutetium is primary, sharing both valence electron count and valence orbital type.
History
The discovery of the group 3 elements is inextricably tied to that of the rare earths, with which they are universally associated in nature. In 1787, Swedish part-time chemist Carl Axel Arrhenius found a heavy black rock near the Swedish village of Ytterby, Sweden (part of the Stockholm Archipelago). Thinking that it was an unknown mineral containing the newly discovered element tungsten, he named it ytterbite. Finnish scientist Johan Gadolin identified a new oxide or "earth" in Arrhenius' sample in 1789, and published his completed analysis in 1794; in 1797, the new oxide was named yttria. In the decades after French scientist Antoine Lavoisier developed the first modern definition of chemical elements, it was believed that earths could be reduced to their elements, meaning that the discovery of a new earth was equivalent to the discovery of the element within, which in this case would have been yttrium. Until the early 1920s, the chemical symbol "Yt" was used for the element, after which "Y" came into common use. Yttrium metal, albeit impure, was first prepared in 1828 when Friedrich Wöhler heated anhydrous yttrium(III) chloride with potassium to form metallic yttrium and potassium chloride. In fact, Gadolin's yttria proved to be a mixture of many metal oxides, that started the history of the discovery of the rare earths.
In 1869, Russian chemist Dmitri Mendeleev published his periodic table, which had an empty space for an element above yttrium. Mendeleev made several predictions on this hypothetical element, which he called eka-boron. By then, Gadolin's yttria had already been split several times; first by Swedish chemist Carl Gustaf Mosander, who in 1843 had split out two more earths which he called terbia and erbia (splitting the name of Ytterby just as yttria had been split); and then in 1878 when Swiss chemist Jean Charles Galissard de Marignac split terbia and erbia themselves into more earths. Among these was ytterbia (a component of the old erbia), which Swedish chemist Lars Fredrik Nilson successfully split in 1879 to reveal yet another new element. He named it scandium, from the Latin Scandia meaning "Scandinavia". Nilson was apparently unaware of Mendeleev's prediction, but Per Teodor Cleve recognized the correspondence and notified Mendeleev. Chemical experiments on scandium proved that Mendeleev's suggestions were correct; along with discovery and characterization of gallium and germanium this proved the correctness of the whole periodic table and periodic law. Metallic scandium was produced for the first time in 1937 by electrolysis of a eutectic mixture, at 700–800 °C, of potassium, lithium, and scandium chlorides. Scandium exists in the same ores that yttrium had been discovered from, but is much rarer and probably for that reason had eluded discovery.
The remaining component of Marignac's ytterbia also proved to be a composite. In 1907, French scientist Georges Urbain, Austrian mineralogist Baron Carl Auer von Welsbach, and American chemist Charles James all independently discovered a new element within ytterbia. Welsbach proposed the name cassiopeium for his new element (after Cassiopeia), whereas Urbain chose the name lutecium (from Latin Lutetia, for Paris). The dispute on the priority of the discovery is documented in two articles in which Urbain and von Welsbach accuse each other of publishing results influenced by the published research of the other. In 1909, the Commission on Atomic Mass, which was responsible for the attribution of the names for the new elements, granted priority to Urbain and adopting his names as official ones. An obvious problem with this decision was that Urbain was one of the four members of the commission. In 1949, the spelling of element 71 was changed to lutetium. Later work connected with Urbain's attempts to further split his lutecium however revealed that it had only contained traces of the new element 71, and that it was only von Welsbach's cassiopeium that was pure element 71. For this reason many German scientists continued to use the name cassiopeium for the element until the 1950s. Ironically, Charles James, who had modestly stayed out of the argument as to priority, worked on a much larger scale than the others, and undoubtedly possessed the largest supply of lutetium at the time. Lutetium was the last of the stable rare earths to be discovered. Over a century of research had split the original yttrium of Gadolin into yttrium, scandium, lutetium, and seven other new elements.
Lawrencium is the only element of the group that does not occur naturally. It was probably first synthesized by Albert Ghiorso and his team on February 14, 1961, at the Lawrence Radiation Laboratory (now called the Lawrence Berkeley National Laboratory) at the University of California in Berkeley, California, United States. The first atoms of lawrencium were produced by bombarding a three-milligram target consisting of three isotopes of the element californium with boron-10 and boron-11 nuclei from the Heavy Ion Linear Accelerator (HILAC). The nuclide 257103 was originally reported. The team at the University of California suggested the name lawrencium (after Ernest O. Lawrence, the inventor of cyclotron particle accelerator) and the symbol "Lw", for the new element; IUPAC accepted their discovery, but changed the symbol to "Lr". In 1965, nuclear-physics researchers in Dubna, Soviet Union (now Russia) reported 256103, in 1967, they reported that they were not able to confirm American scientists' data on 257103, and proposed the name "rutherfordium" for the new element. The Dubna group criticised the IUPAC approval of the Berkeley group's discovery as having been hasty. In 1971, the Berkeley group did a whole series of experiments aimed at measuring the nuclear decay properties of element 103 isotopes, in which all previous results from Berkeley and Dubna were confirmed, except that the initial 257103 isotope reported at Berkeley in 1961 turned out to have been 258103. In 1992, the IUPAC Trans-fermium Working Group named the nuclear physics teams at Dubna and Berkeley as the co-discoverers of element 103. When IUPAC made the final decision of the naming of the elements beyond 100 in 1997, it decided to keep the name "lawrencium" and symbol "Lr" for element 103 as it had been in use for a long time by that point. The name "rutherfordium" was assigned to the following element 104, which the Berkeley team had proposed it for.
Characteristics
Chemical
Like other groups, the members of this family show patterns in their electron configurations, especially the outermost shells, resulting in trends in chemical behavior. Due to relativistic effects that become important for high atomic numbers, lawrencium's configuration has an irregular 7p occupancy instead of the expected 6d, but the regular [Rn]5f146d17s2 configuration is low enough in energy that no significant difference from the rest of the group is observed or expected.
Most of the chemistry has been observed only for the first three members of the group; chemical properties of lawrencium are not well-characterized, but what is known and predicted matches its position as a heavier homolog of lutetium. The remaining elements of the group (scandium, yttrium, lutetium) are quite electropositive. They are reactive metals, although this is not obvious due to the formation of a stable oxide layer which prevents further reactions. The metals burn easily to give the oxides, which are white high-melting solids. They are usually oxidized to the +3 oxidation state, in which they form mostly ionic compounds and have a mostly cationic aqueous chemistry. In this way they are similar to the lanthanides, although they lack the involvement of f orbitals that characterises the chemistry of the 4f elements lanthanum through ytterbium. The stable group 3 elements are thus often grouped with the 4f elements as the so-called rare earths.
Typical transition-metal properties are mostly absent from this group, as they are for the heavier elements of groups 4 and 5: there is only one typical oxidation state and the coordination chemistry is not very rich (though high coordination numbers are common due to the large size of the M3+ ions). This said, low-oxidation state compounds may be prepared and some cyclopentadienyl chemistry is known. The chemistries of group 3 elements are thus mostly distinguished by their atomic radii: yttrium and lutetium are very similar, but scandium stands out as the least basic and the best complexing agent, approaching aluminium in some properties. They naturally take their places together with the rare earths in a series of trivalent elements: yttrium acts as a rare earth intermediate between dysprosium and holmium in basicity; lutetium as less basic than the 4f elements and the least basic of the lanthanides; and scandium as a rare earth less basic than even lutetium. Scandium oxide is amphoteric; lutetium oxide is more basic (although it can with difficulty be made to display some acidic properties), and yttrium oxide is more basic still. Salts with strong acids of these metals are soluble, whereas those with weak acids (e.g. fluorides, phosphates, oxalates) are sparingly soluble or insoluble.
Physical
The trends in group 3 follow those of the other early d-block groups and reflect the addition of a filled f-shell into the core in passing from the fifth to the sixth period. For example, scandium and yttrium are both soft metals. But because of the lanthanide contraction, the expected increase in atomic radius from yttrium to lutetium is reversed; lutetium atoms are slightly smaller than yttrium atoms, but are heavier and have a higher nuclear charge. This makes the metal more dense, and also harder because the extraction of the electrons from the atom to form metallic bonding becomes more difficult. All three metals have similar melting and boiling points. Very little is known about lawrencium, but calculations suggest it continues the trend of its lighter congeners toward increasing density.
Scandium, yttrium, and lutetium all crystallize in the hexagonal close-packed structure at room temperature, and lawrencium is expected to do the same. The stable members of the group are known to change structure at high temperature. In comparison with most metals, they are not very good conductors of heat and electricity because of the low number of electrons available for metallic bonding.
Occurrence
Scandium, yttrium, and lutetium tend to occur together with the other lanthanides (except short-lived promethium) in the Earth's crust, and are often harder to extract from their ores. The abundance of elements in Earth's crust for group 3 is quite low—all the elements in the group are uncommon, the most abundant being yttrium with abundance of approximately 30 parts per million (ppm); the abundance of scandium is 16 ppm, while that of lutetium is about 0.5 ppm. For comparison, the abundance of copper is 50 ppm, that of chromium is 160 ppm, and that of molybdenum is 1.5 ppm.
Scandium is distributed sparsely and occurs in trace amounts in many minerals. Rare minerals from Scandinavia and Madagascar such as gadolinite, euxenite, and thortveitite are the only known concentrated sources of this element, the latter containing up to 45% of scandium in the form of scandium(III) oxide. Yttrium has the same trend in occurrence places; it is found in lunar rock samples collected during the American Apollo Project in a relatively high content as well.
The principal commercially viable ore of lutetium is the rare-earth phosphate mineral monazite, (Ce,La,etc.)PO4, which contains 0.003% of the element. The main mining areas are China, United States, Brazil, India, Sri Lanka and Australia. Pure lutetium metal is one of the rarest and most expensive of the rare-earth metals with the price about US$10,000/kg, or about one-fourth that of gold.
Production
The most available element in group 3 is yttrium, with annual production of 8,900 tonnes in 2010. Yttrium is mostly produced as oxide, by a single country, China (99%). Lutetium and scandium are also mostly obtained as oxides, and their annual production by 2001 was about 10 and 2 tonnes, respectively.
Group 3 elements are mined only as a byproduct from the extraction of other elements. They are not often produced as the pure metals; the production of metallic yttrium is about a few tonnes, and that of scandium is in the order of 10 kg per year; production of lutetium is not calculated, but it is certainly small. The elements, after purification from other rare-earth metals, are isolated as oxides; the oxides are converted to fluorides during reactions with hydrofluoric acid. The resulting fluorides are reduced with alkaline earth metals or alloys of the metals; metallic calcium is used most frequently. For example:
Sc2O3 + 3 HF → 2 ScF3 + 3 H2O
2 ScF3 + 3 Ca → 3 CaF2 + 2 Sc
Biological chemistry
Group 3 metals have low availability to the biosphere. Scandium, yttrium, and lutetium have no documented biological role in living organisms. The high radioactivity of lawrencium would make it highly toxic to living cells, causing radiation poisoning.
Scandium concentrates in the liver and is a threat to it; some of its compounds are possibly carcinogenic, even though in general scandium is not toxic. Scandium is known to have reached the food chain, but in trace amounts only; a typical human takes in less than 0.1 micrograms per day. Once released into the environment, scandium gradually accumulates in soils, which leads to increased concentrations in soil particles, animals and humans. Scandium is mostly dangerous in the working environment, due to the fact that damps and gases can be inhaled with air. This can cause lung embolisms, especially during long-term exposure. The element is known to damage cell membranes of water animals, causing several negative influences on reproduction and on the functions of the nervous system.
Yttrium tends to concentrate in the liver, kidney, spleen, lungs, and bones of humans. There is normally as little as 0.5 milligrams found within the entire human body; human breast milk contains 4 ppm. Yttrium can be found in edible plants in concentrations between 20 ppm and 100 ppm (fresh weight), with cabbage having the largest amount. With up to 700 ppm, the seeds of woody plants have the highest known concentrations.
Lutetium concentrates in bones, and to a lesser extent in the liver and kidneys. Lutetium salts are known to cause metabolism and they occur together with other lanthanide salts in nature; the element is the least abundant in the human body of all lanthanides. Human diets have not been monitored for lutetium content, so it is not known how much the average human takes in, but estimations show the amount is only about several micrograms per year, all coming from tiny amounts taken by plants. Soluble lutetium salts are mildly toxic, but insoluble ones are not.
| Physical sciences | Group 3 | Chemistry |
306755 | https://en.wikipedia.org/wiki/Spitzer%20Space%20Telescope | Spitzer Space Telescope | The Spitzer Space Telescope, formerly the Space Infrared Telescope Facility (SIRTF), was an infrared space telescope launched in 2003, that was deactivated when operations ended on 30 January 2020. Spitzer was the third space telescope dedicated to infrared astronomy, following IRAS (1983) and ISO (1995–1998). It was the first spacecraft to use an Earth-trailing orbit, later used by the Kepler planet-finder.
The planned mission period was to be 2.5 years with a pre-launch expectation that the mission could extend to five or slightly more years until the onboard liquid helium supply was exhausted. This occurred on 15 May 2009. Without liquid helium to cool the telescope to the very low temperatures needed to operate, most of the instruments were no longer usable. However, the two shortest-wavelength modules of the IRAC camera continued to operate with the same sensitivity as before the helium was exhausted, and continued to be used into early 2020 in the Spitzer Warm Mission.
During the warm mission, the two short wavelength channels of IRAC operated at 28.7 K and were predicted to experience little to no degradation at this temperature compared to the nominal mission. The Spitzer data, from both the primary and warm phases, are archived at the Infrared Science Archive (IRSA).
In keeping with NASA tradition, the telescope was renamed after its successful demonstration of operation, on 18 December 2003. Unlike most telescopes that are named by a board of scientists, typically after famous deceased astronomers, the new name for SIRTF was obtained from a contest open to the general public.
The contest led to the telescope being named in honor of astronomer Lyman Spitzer, who had promoted the concept of space telescopes in the 1940s. Spitzer wrote a 1946 report for RAND Corporation describing the advantages of an extraterrestrial observatory and how it could be realized with available or upcoming technology. He has been cited for his pioneering contributions to rocketry and astronomy, as well as "his vision and leadership in articulating the advantages and benefits to be realized from the Space Telescope Program."
The Spitzer was launched on 25 August 2003 at 05:35:39 UTC from Cape Canaveral SLC-17B aboard a Delta II 7920H rocket. It was placed into a heliocentric (as opposed to a geocentric) orbit trailing and drifting away from Earth's orbit at approximately 0.1 astronomical units per year (an "Earth-trailing" orbit).
The primary mirror is in diameter, , made of beryllium and was cooled to . The satellite contains three instruments that allowed it to perform astronomical imaging and photometry from 3.6 to 160 micrometers, spectroscopy from 5.2 to 38 micrometers, and spectrophotometry from 55 to 95 micrometers.
History
By the early 1970s, astronomers began to consider the possibility of placing an infrared telescope above the obscuring effects of Earth's atmosphere.
In 1979, a report from the National Research Council of the National Academy of Sciences, A Strategy for Space Astronomy and Astrophysics for the 1980s, identified a Shuttle Infrared Telescope Facility (SIRTF) as "one of two major astrophysics facilities [to be developed] for Spacelab", a shuttle-borne platform. Anticipating the major results from an upcoming Explorer satellite and from the Shuttle mission, the report also favored the "study and development of ... long-duration spaceflights of infrared telescopes cooled to cryogenic temperatures."
The launch in January 1983 of the Infrared Astronomical Satellite, jointly developed by the United States, the Netherlands, and the United Kingdom, to conduct the first infrared survey of the sky, whetted the appetites of scientists worldwide for follow-up space missions capitalizing on the rapid improvements in infrared detector technology.
Earlier infrared observations had been made by both space-based and ground-based observatories. Ground-based observatories have the drawback that at infrared wavelengths or frequencies, both the Earth's atmosphere and the telescope itself will radiate (glow) brightly. Additionally, the atmosphere is opaque at most infrared wavelengths. This necessitates lengthy exposure times and greatly decreases the ability to detect faint objects. It could be compared to trying to observe the stars in the optical at noon from a telescope built out of light bulbs. Previous space observatories (such as IRAS, the Infrared Astronomical Satellite, and ISO, the Infrared Space Observatory) were launched during the 1980s and 1990s and great advances in astronomical technology have been made since then.
Most of the early concepts envisioned repeated flights aboard the NASA Space Shuttle. This approach was developed in an era when the Shuttle program was expected to support weekly flights of up to 30 days duration. A May 1983 NASA proposal described SIRTF as a Shuttle-attached mission, with an evolving scientific instrument payload. Several flights were anticipated with a probable transition into a more extended mode of operation, possibly in association with a future space platform or space station. SIRTF would be a 1-meter class, cryogenically cooled, multi-user facility consisting of a telescope and associated focal plane instruments. It would be launched on the Space Shuttle and remain attached to the Shuttle as a Spacelab payload during astronomical observations, after which it would be returned to Earth for refurbishment prior to re-flight. The first flight was expected to occur about 1990, with the succeeding flights anticipated beginning approximately one year later. However, the Spacelab-2 flight aboard STS-51-F showed that the Shuttle environment was poorly suited to an onboard infrared telescope due to contamination from the relatively "dirty" vacuum associated with the orbiters. By September 1983, NASA was considering the "possibility of a long duration [free-flyer] SIRTF mission".
Spitzer is the only one of the Great Observatories not launched by the Space Shuttle, as was originally intended. However, after the 1986 Challenger disaster, the Shuttle-Centaur upper stage, which would have been required to place it into its final orbit, was abandoned. The mission underwent a series of redesigns during the 1990s, primarily due to budget considerations. This resulted in a much smaller but still fully capable mission that could use the smaller Delta II expendable launch vehicle.
One of the most important advances of this redesign was an Earth-trailing orbit. Cryogenic satellites that require liquid helium (LHe, T ≈ 4 K) temperatures in near-Earth orbit are typically exposed to a large heat load from Earth, and consequently require large amounts of LHe coolant, which then tends to dominate the total payload mass and limits mission life. Placing the satellite in solar orbit far from Earth allowed innovative passive cooling. The sun shield protected the rest of the spacecraft from the Sun's heat, the far side of the spacecraft was painted black to enhance passive radiation of heat, and the spacecraft bus was thermally isolated from the telescope. All of these design choices combined to drastically reduce the total mass of helium needed, resulting in an overall smaller and lighter payload, resulting in major cost savings, but with a mirror the same diameter as originally designed. This orbit also simplified telescope pointing, but did require the NASA Deep Space Network for communications.
The primary instrument package (telescope and cryogenic chamber) was developed by Ball Aerospace & Technologies, in Boulder, Colorado. The individual instruments were developed jointly by industrial, academic, and government institutions. The principals were Cornell University, the University of Arizona, the Smithsonian Astrophysical Observatory, Ball Aerospace, and Goddard Spaceflight Center. The shorter-wavelength infrared detectors were developed by Raytheon in Goleta, California. Raytheon used indium antimonide and a doped silicon detector in the creation of the infrared detectors. These detectors are 100 times more sensitive than what was available at the beginning of the project during the 1980s. The far-infrared detectors (70–160 micrometers) were developed jointly by the University of Arizona and Lawrence Berkeley National Laboratory using gallium-doped germanium. The spacecraft was built by Lockheed Martin. The mission was operated and managed by the Jet Propulsion Laboratory and the Spitzer Science Center, located at IPAC on the Caltech campus in Pasadena, California.
Launch and commissioning
Warm mission and end of mission
Spitzer ran out of liquid helium coolant on 15 May 2009, which stopped far-IR observations. Only the IRAC instrument remained in use, and only at the two shorter wavelength bands (3.6 μm and 4.5 μm). The telescope equilibrium temperature was then around , and IRAC continued to produce valuable images at those wavelengths as the "Spitzer Warm Mission".
Late in the mission, ~2016, Spitzer's distance to Earth and the shape of its orbit meant the spacecraft had to pitch over at an extreme angle to aim its antenna at Earth. The solar panels were not fully illuminated at this angle, and this limited those communications to 2.5 hours due to the battery drain. The telescope was retired on 30 January 2020 when NASA sent a shutdown signal to the telescope from the Goldstone Deep Space Communications Complex (GDSCC) instructing the telescope to go into safe mode. After receiving confirmation that the command was successful, Spitzer Project Manager Joseph Hunt officially declared that the mission had ended.
Instruments
Spitzer carries three instruments on board:
Infrared Array Camera (IRAC)
An infrared camera which operated simultaneously on four wavelengths (3.6 μm, 4.5 μm, 5.8 μm and 8 μm). Each module used a 256×256-pixel detector—the short-wavelength pair used indium antimonide technology, the long-wavelength pair used arsenic-doped silicon impurity band conduction technology. The principal investigator was Giovanni Fazio of Center for Astrophysics Harvard & Smithsonian; the flight hardware was built by NASA Goddard Space Flight Center.
Infrared Spectrograph (IRS)
An infrared spectrometer with four sub-modules which operate at the wavelengths 5.3–14 μm (low resolution), 10–19.5 μm (high resolution), 14–40 μm (low resolution), and 19–37 μm (high resolution). Each module used a 128×128-pixel detector—the short-wavelength pair used arsenic-doped silicon blocked impurity band technology, the long-wavelength pair used antimony-doped silicon blocked impurity band technology. The principal investigator was James R. Houck of Cornell University; the flight hardware was built by Ball Aerospace.
Multiband Imaging Photometer for Spitzer (MIPS)
Three detector arrays in the mid- to far-infrared (128 × 128 pixels at 24 μm, 32 × 32 pixels at 70 μm, 2 × 20 pixels at 160 μm). The 24 μm detector is identical to one of the IRS short-wavelength modules. The 70 μm detector used gallium-doped germanium technology, and the 160 μm detector also used gallium-doped germanium, but with mechanical stress added to each pixel to lower the bandgap and extend sensitivity to this long-wavelength. The principal investigator was George H. Rieke of the University of Arizona; the flight hardware was built by Ball Aerospace.
All three instruments used liquid helium for cooling the sensors. Once the helium was exhausted, only the two shorter wavelengths in IRAC were used in the "warm mission".
Results
While some time on the telescope was reserved for participating institutions and crucial projects, astronomers around the world also had the opportunity to submit proposals for observing time. Prior to launch, there was a proposal call for large, coherent investigations using Spitzer. If the telescope failed early and/or ran out of cryogen very quickly, these so-called Legacy Projects would ensure the best possible science could be obtained quickly in the early months of the mission. As a requirement tied to the funding these Legacy teams received, the teams had to deliver high-level data products back to the Spitzer Science Center (and the NASA/IPAC Infrared Science Archive) for use by the community, again ensuring the rapid scientific return of the mission. The international scientific community quickly realized the value of delivering products for others to use, and even though Legacy projects were no longer explicitly solicited in subsequent proposal calls, teams continued to deliver products to the community. The Spitzer Science Center later reinstated named "Legacy" projects (and later still "Exploration Science" projects) in response to this community-driven effort.
Important targets included forming stars (young stellar objects, or YSOs), planets, and other galaxies. Images are freely available for educational and journalistic purposes.
The first released images from Spitzer were designed to show off the abilities of the telescope and showed a glowing stellar nursery, a big swirling, dusty galaxy, a disc of planet-forming debris, and organic material in the distant universe. Since then, many monthly press releases have highlighted Spitzer capabilities, as the NASA and ESA images do for the Hubble Space Telescope.
As one of its most noteworthy observations, in 2005, Spitzer became one of the first telescopes to directly capture light from exoplanets, namely the "hot Jupiters" HD 209458 b and TrES-1b, although it did not resolve that light into actual images. This was one of the first times the light from extrasolar planets had been directly detected; earlier observations had been indirectly made by drawing conclusions from behaviors of the stars the planets were orbiting. The telescope also discovered in April 2005 that Cohen-kuhi Tau/4 had a planetary disk that was vastly younger and contained less mass than previously theorized, leading to new understandings of how planets are formed.
In 2004, it was reported that Spitzer had spotted a faintly glowing body that may be the youngest star ever seen. The telescope was trained on a core of gas and dust known as L1014 which had previously appeared completely dark to ground-based observatories and to ISO (Infrared Space Observatory), a predecessor to Spitzer. The advanced technology of Spitzer revealed a bright red hot spot in the middle of L1014.
Scientists from the University of Texas at Austin, who discovered the object, believe the hot spot to be an example of early star development, with the young star collecting gas and dust from the cloud around it. Early speculation about the hot spot was that it might have been the faint light of another core that lies 10 times further from Earth but along the same line of sight as L1014. Follow-up observation from ground-based near-infrared observatories detected a faint fan-shaped glow in the same location as the object found by Spitzer. That glow is too feeble to have come from the more distant core, leading to the conclusion that the object is located within L1014. (Young et al., 2004)
In 2005, astronomers from the University of Wisconsin at Madison and Whitewater determined, on the basis of 400 hours of observation on the Spitzer Space Telescope, that the Milky Way galaxy has a more substantial bar structure across its core than previously recognized.
Also in 2005, astronomers Alexander Kashlinsky and John Mather of NASA's Goddard Space Flight Center reported that one of Spitzer earliest images may have captured the light of the first stars in the universe. An image of a quasar in the Draco constellation, intended only to help calibrate the telescope, was found to contain an infrared glow after the light of known objects was removed. Kashlinsky and Mather are convinced that the numerous blobs in this glow are the light of stars that formed as early as 100 million years after the Big Bang, redshifted by cosmic expansion.
In March 2006, astronomers reported an nebula near the center of the Milky Way Galaxy, the Double Helix Nebula, which is, as the name implies, twisted into a double spiral shape. This is thought to be evidence of massive magnetic fields generated by the gas disc orbiting the supermassive black hole at the galaxy's center, from the nebula and from Earth. This nebula was discovered by Spitzer and published in the magazine Nature on 16 March 2006.
In May 2007, astronomers successfully mapped the atmospheric temperature of HD 189733 b, thus obtaining the first map of some kind of an extrasolar planet.
Starting in September 2006, the telescope participated in a series of surveys called the Gould Belt Survey, observing the Gould's Belt region in multiple wavelengths. The first set of observations by the Spitzer Space Telescope was completed from 21 September 2006 through 27 September. Resulting from these observations, the team of astronomers led by Dr. Robert Gutermuth, of the Center for Astrophysics Harvard & Smithsonian reported the discovery of Serpens South, a cluster of 50 young stars in the Serpens constellation.
Scientists have long wondered how tiny silicate crystals, which need high temperatures to form, have found their way into frozen comets, born in the very cold environment of the Solar System's outer edges. The crystals would have begun as non-crystallized, amorphous silicate particles, part of the mix of gas and dust from which the Solar System developed. This mystery has deepened with the results of the Stardust sample return mission, which captured particles from Comet Wild 2. Many of the Stardust particles were found to have formed at temperatures in excess of 1,000 K.
In May 2009, Spitzer researchers from Germany, Hungary, and the Netherlands found that amorphous silicate appears to have been transformed into crystalline form by an outburst from a star. They detected the infrared signature of forsterite silicate crystals on the disk of dust and gas surrounding the star EX Lupi during one of its frequent flare-ups, or outbursts, seen by Spitzer in April 2008. These crystals were not present in Spitzer previous observations of the star's disk during one of its quiet periods. These crystals appear to have formed by radiative heating of the dust within 0.5 AU of EX Lupi.
In August 2009, the telescope found evidence of a high-speed collision between two burgeoning planets orbiting a young star.
In October 2009, astronomers Anne J. Verbiscer, Michael F. Skrutskie, and Douglas P. Hamilton published findings of the "Phoebe ring" of Saturn, which was found with the telescope; the ring is a huge, tenuous disc of material extending from 128 to 207 times the radius of Saturn.
GLIMPSE and MIPSGAL surveys
GLIMPSE, the Galactic Legacy Infrared Mid-Plane Survey Extraordinaire, was a series of surveys that spanned 360° of the inner region of the Milky Way galaxy, which provided the first large-scale mapping of the galaxy. It consists of more than 2 million snapshots taken in four separate wavelengths using the Infrared Array Camera. The images were taken over a 10-year period beginning in 2003 when Spitzer launched.
MIPSGAL, a similar survey that complements GLIMPSE, covers 248° of the galactic disk using the 24 and 70 μm channels of the MIPS instrument.
On 3 June 2008, scientists unveiled the largest, most detailed infrared portrait of the Milky Way, created by stitching together more than 800,000 snapshots, at the 212th meeting of the American Astronomical Society in St. Louis, Missouri. This composite survey is now viewable with the GLIMPSE/MIPSGAL Viewer.
2010s
Spitzer observations, announced in May 2011, indicate that tiny forsterite crystals might be falling down like rain on to the protostar HOPS-68. The discovery of the forsterite crystals in the outer collapsing cloud of the protostar is surprising because the crystals form at lava-like high temperatures, yet they are found in the molecular cloud where the temperatures are about . This led the team of astronomers to speculate that the bipolar outflow from the young star may be transporting the forsterite crystals from near the star's surface to the chilly outer cloud.
In January 2012, it was reported that further analysis of the Spitzer observations of EX Lupi can be understood if the forsterite crystalline dust was moving away from the protostar at a remarkable average speed of . It would appear that such high speeds can arise only if the dust grains had been ejected by a bipolar outflow close to the star. Such observations are consistent with an astrophysical theory, developed in the early 1990s, where it was suggested that bipolar outflows garden or transform the disks of gas and dust that surround protostars by continually ejecting reprocessed, highly heated material from the inner disk, adjacent to the protostar, to regions of the accretion disk further away from the protostar.
In April 2015, Spitzer and the Optical Gravitational Lensing Experiment were reported as co-discovering one of the most distant planets ever identified: a gas giant about away from Earth.
In June and July 2015, the brown dwarf was discovered using the gravitational microlensing detection method in a joint effort between Swift, Spitzer, and the ground-based Optical Gravitational Lensing Experiment, the first time two space telescopes have observed the same microlensing event. This method was possible because of the large separation between the two spacecraft: Swift is in low-Earth orbit while Spitzer is more than one AU distant in an Earth-trailing heliocentric orbit. This separation provided significantly different perspectives of the brown dwarf, allowing for constraints to be placed on some of the object's physical characteristics.
Reported in March 2016, Spitzer and Hubble were used to discover the most distant-known galaxy, GN-z11. This object was seen as it appeared 13.4 billion years ago.
Spitzer Beyond
On 1 October 2016, Spitzer began its Observation Cycle 13, a year extended mission nicknamed Beyond. One of the goals of this extended mission was to help prepare for the James Webb Space Telescope, also an infrared telescope, by identifying candidates for more detailed observations.
Another aspect of the Beyond mission was the engineering challenges of operating Spitzer in its progressing orbital phase. As the spacecraft moved farther from Earth on the same orbital path from the Sun, its antenna had to point at increasingly higher angles to communicate with ground stations; this change in angle imparted more and more solar heating on the vehicle while its solar panels received less sunlight.
Planet hunter
Spitzer was also put to work studying exoplanets thanks to creatively tweaking its hardware. This included doubling its stability by modifying its heating cycle, finding a new use for the "peak-up" camera, and analyzing the sensor at a sub-pixel level. Although in its "warm" mission, the spacecraft's passive cooling system kept the sensors at . Spitzer used the transit photometry and gravitational microlensing techniques to perform these observations. According to NASA's Sean Carey, "We never even considered using Spitzer for studying exoplanets when it launched. ... It would have seemed ludicrous back then, but now it's an important part of what Spitzer does."
Examples of exoplanets discovered using Spitzer include HD 219134 b in 2015, which was shown to be a rocky planet about 1.5 times as large as Earth in a three-day orbit around its star; and an unnamed planet found using microlensing located about from Earth.
In September–October 2016, Spitzer was used to discover five of a total of seven known planets around the star TRAPPIST-1, all of which are approximately Earth-sized and likely rocky. Three of the discovered planets are located in the habitable zone, which means they are capable of supporting liquid water given sufficient parameters. Using the transit method, Spitzer helped measure the sizes of the seven planets and estimate the mass and density of the inner six. Further observations will help determine if there is liquid water on any of the planets.
| Technology | Space-based observatories | null |
306773 | https://en.wikipedia.org/wiki/Water%20quality | Water quality | Water quality refers to the chemical, physical, and biological characteristics of water based on the standards of its usage. It is most frequently used by reference to a set of standards against which compliance, generally achieved through treatment of the water, can be assessed. The most common standards used to monitor and assess water quality convey the health of ecosystems, safety of human contact, extent of water pollution and condition of drinking water. Water quality has a significant impact on water supply and often determines supply options.
Impacts on public health
Over time, there has been increasing recognition of the importance of drinking water quality and its impact on public health. This has led to increasing protection and management of water quality.
The understanding of the links between water quality and health continues to grow and highlight new potential health crises: from the chronic impacts of infectious diseases on child development through stunting to new evidence on the harms from known contaminants, such as manganese with growing evidence of neurotoxicity in children. In addition, there are many emerging water quality issues—such as microplastics, perfluorinated compounds, and antimicrobial resistance.
Categories
The parameters for water quality are determined by the intended use. Work in the area of water quality tends to be focused on water that is treated for potability, industrial/domestic use, or restoration (of an environment/ecosystem, generally for health of human/aquatic life).
Human consumption
Contaminants that may be in untreated water include microorganisms such as viruses, protozoa and bacteria; inorganic contaminants such as salts and metals; organic chemical contaminants from industrial processes and petroleum use; pesticides and herbicides; and radioactive contaminants. Water quality depends on the local geology and ecosystem, as well as human uses such as sewage dispersion, industrial pollution, use of water bodies as a heat sink, and overuse (which may lower the level of the water).
The United States Environmental Protection Agency (EPA) limits the amounts of certain contaminants in tap water provided by US public water systems. The Safe Drinking Water Act authorizes EPA to issue two types of standards:
primary standards regulate substances that potentially affect human health;
secondary standards prescribe aesthetic qualities, those that affect taste, odor, or appearance.
The U.S. Food and Drug Administration (FDA) regulations establish limits for contaminants in bottled water. Drinking water, including bottled water, may reasonably be expected to contain at least small amounts of some contaminants. The presence of these contaminants does not necessarily indicate that the water poses a health risk.
In urbanized areas around the world, water purification technology is used in municipal water systems to remove contaminants from the source water (surface water or groundwater) before it is distributed to homes, businesses, schools and other recipients. Water drawn directly from a stream, lake, or aquifer and that has no treatment will be of uncertain quality in terms of potability.
The burden of polluted drinking water disproportionally effects under-represented and vulnerable populations. Communities that lack these clean drinking-water services are at risk of contracting water-borne and pollution-related illnesses like Cholera, diarrhea, dysentery, hepatitis A, typhoid, and polio. These communities are often in low-income areas, where human wastewater is discharged into a nearby drainage channel or surface water drain without sufficient treatment, or is used in agricultural irrigation.
Industrial and domestic use
Dissolved ions may affect the suitability of water for a range of industrial and domestic purposes. The most familiar of these is probably the presence of calcium (Ca2+) and magnesium (Mg2+) that interfere with the cleaning action of soap, and can form hard sulfate and soft carbonate deposits in water heaters or boilers. Hard water may be softened to remove these ions. The softening process often substitutes sodium cations. For certain populations, hard water may be preferable to soft water because health problems have been associated with calcium deficiencies and with excess sodium. The necessity for additional calcium and magnesium in water depends on the population in question because people generally satisfy their recommended amounts through food.
Environmental water quality
Environmental water quality, also called ambient water quality, relates to water bodies such as lakes, rivers, and oceans. Water quality standards for surface waters vary significantly due to different environmental conditions, ecosystems, and intended human uses. Toxic substances and high populations of certain microorganisms can present a health hazard for non-drinking purposes such as irrigation, swimming, fishing, rafting, boating, and industrial uses. These conditions may also affect wildlife, which use the water for drinking or as a habitat. According to the EPA, water quality laws generally specify protection of fisheries and recreational use and require, as a minimum, retention of current quality standards. In some locations, desired water quality conditions include high dissolved oxygen concentrations, low chlorophyll-a concentrations, and high water clarity.
There is some desire among the public to return water bodies to pristine, or pre-industrial conditions. Most current environmental laws focus on the designation of particular uses of a water body. In some countries these designations allow for some water contamination as long as the particular type of contamination is not harmful to the designated uses. Given the landscape changes (e.g., land development, urbanization, clearcutting in forested areas) in the watersheds of many freshwater bodies, returning to pristine conditions would be a significant challenge. In these cases, environmental scientists focus on achieving goals for maintaining healthy ecosystems and may concentrate on the protection of populations of endangered species and protecting human health.
Sampling and measurement
Sample collection
The complexity of water quality as a subject is reflected in the many types of measurements of water quality indicators. Some measurements of water quality are most accurately made on-site, because water exists in equilibrium with its surroundings. Measurements commonly made on-site and in direct contact with the water source in question include temperature, pH, dissolved oxygen, conductivity, oxygen reduction potential (ORP), turbidity, and Secchi disk depth.
Sampling of water for physical or chemical testing can be done by several methods, depending on the accuracy needed and the characteristics of the contaminant. Sampling methods include for example simple random sampling, stratified sampling, systematic and grid sampling, adaptive cluster sampling, grab samples, semi-continuous monitoring and continuous, passive sampling, remote surveillance, remote sensing, and biomonitoring. The use of passive samplers greatly reduces the cost and the need of infrastructure on the sampling location.
Many contamination events are sharply restricted in time, most commonly in association with rain events. For this reason "grab" samples are often inadequate for fully quantifying contaminant levels. Scientists gathering this type of data often employ auto-sampler devices that pump increments of water at either time or discharge intervals.
More complex measurements are often made in a laboratory requiring a water sample to be collected, preserved, transported, and analyzed at another location.
Issues
The process of water sampling introduces two significant problems:
The first problem is the extent to which the sample may be representative of the water source of interest. Water sources vary with time and with location. The measurement of interest may vary seasonally or from day to night or in response to some activity of man or natural populations of aquatic plants and animals. The measurement of interest may vary with distances from the water boundary with overlying atmosphere and underlying or confining soil. The sampler must determine if a single time and location meets the needs of the investigation, or if the water use of interest can be satisfactorily assessed by averaged values of sampling over time and location, or if critical maxima and minima require individual measurements over a range of times, locations or events. The sample collection procedure must assure correct weighting of individual sampling times and locations where averaging is appropriate. Where critical maximum or minimum values exist, statistical methods must be applied to observed variation to determine an adequate number of samples to assess the probability of exceeding those critical values.
The second problem occurs as the sample is removed from the water source and begins to establish chemical equilibrium with its new surroundings – the sample container. Sample containers must be made of materials with minimal reactivity with substances to be measured; pre-cleaning of sample containers is important. The water sample may dissolve part of the sample container and any residue on that container, and chemicals dissolved in the water sample may sorb onto the sample container and remain there when the water is poured out for analysis. Similar physical and chemical interactions may take place with any pumps, piping, or intermediate devices used to transfer the water sample into the sample container. Water collected from depths below the surface will normally be held at the reduced pressure of the atmosphere; so gas dissolved in the water will collect at the top of the container. Atmospheric gas above the water may also dissolve into the water sample. Other chemical reaction equilibria may change if the water sample changes temperature. Finely divided solid particles formerly suspended by water turbulence may settle to the bottom of the sample container, or a solid phase may form from biological growth or chemical precipitation. Microorganisms within the water sample may biochemically alter concentrations of oxygen, carbon dioxide, and organic compounds. Changing carbon dioxide concentrations may alter pH and change solubility of chemicals of interest. These problems are of special concern during measurement of chemicals assumed to be significant at very low concentrations.
Sample preservation may partially resolve the second problem. A common procedure is keeping samples cold to slow the rate of chemical reactions and phase change, and analyzing the sample as soon as possible; but this merely minimizes the changes rather than preventing them. A useful procedure for determining influence of sample containers during delay between sample collection and analysis involves preparation for two artificial samples in advance of the sampling event. One sample container is filled with water known from previous analysis to contain no detectable amount of the chemical of interest. This sample, called a "blank", is opened for exposure to the atmosphere when the sample of interest is collected, then resealed and transported to the laboratory with the sample for analysis to determine if sample collection or holding procedures introduced any measurable amount of the chemical of interest. The second artificial sample is collected with the sample of interest, but then "spiked" with a measured additional amount of the chemical of interest at the time of collection. The blank (negative control) and spiked sample (positive control) are carried with the sample of interest and analyzed by the same methods at the same times to determine any changes indicating gains or losses during the elapsed time between collection and analysis.
Testing in response to natural disasters and other emergencies
After events such as earthquakes and tsunamis, there is an immediate response by the aid agencies as relief operations get underway to try and restore basic infrastructure and provide the basic fundamental items that are necessary for survival and subsequent recovery. The threat of disease increases hugely due to the large numbers of people living close together, often in squalid conditions, and without proper sanitation.
After a natural disaster, as far as water quality testing is concerned, there are widespread views on the best course of action to take and a variety of methods can be employed. The key basic water quality parameters that need to be addressed in an emergency are bacteriological indicators of fecal contamination, free chlorine residual, pH, turbidity and possibly conductivity/total dissolved solids. There are many decontamination methods.
After major natural disasters, a considerable length of time might pass before water quality returns to pre-disaster levels. For example, following the 2004 Indian Ocean tsunami the Colombo-based International Water Management Institute (IWMI) monitored the effects of saltwater and concluded that the wells recovered to pre-tsunami drinking water quality one and a half years after the event. IWMI developed protocols for cleaning wells contaminated by saltwater; these were subsequently officially endorsed by the World Health Organization as part of its series of Emergency Guidelines.
Chemical analysis
The simplest methods of chemical analysis are those measuring chemical elements without respect to their form. Elemental analysis for oxygen, as an example, would indicate a concentration of 890 g/L (grams per litre) of water sample because oxygen (O) has 89% mass of the water molecule (H2O). The method selected to measure dissolved oxygen should differentiate between diatomic oxygen and oxygen combined with other elements. The comparative simplicity of elemental analysis has produced a large amount of sample data and water quality criteria for elements sometimes identified as heavy metals. Water analysis for heavy metals must consider soil particles suspended in the water sample. These suspended soil particles may contain measurable amounts of metal. Although the particles are not dissolved in the water, they may be consumed by people drinking the water. Adding acid to a water sample to prevent loss of dissolved metals onto the sample container may dissolve more metals from suspended soil particles. Filtration of soil particles from the water sample before acid addition, however, may cause loss of dissolved metals onto the filter. The complexities of differentiating similar organic molecules are even more challenging.
Making these complex measurements can be expensive. Because direct measurements of water quality can be expensive, ongoing monitoring programs are typically conducted and results released by government agencies. However, there are local volunteer programs and resources available for some general assessment. Tools available to the general public include on-site test kits, commonly used for home fish tanks, and biological assessment procedures.
Biosensors
Biosensors have the potential for "high sensitivity, selectivity, reliability, simplicity, low-cost and real-time response". For instance, bionanotechnologists reported the development of , that can detect levels of diverse water pollutants.
Real-time monitoring
Although water quality is usually sampled and analyzed at laboratories, since the late 20th century there has been increasing public interest in the quality of drinking water provided by municipal systems. Many water utilities have developed systems to collect real-time data about source water quality. In the early 21st century, a variety of sensors and remote monitoring systems have been deployed for measuring water pH, turbidity, dissolved oxygen and other parameters. Some remote sensing systems have also been developed for monitoring ambient water quality in riverine, estuarine and coastal water bodies.
The following is a list of indicators often measured by situational category:
Alkalinity
Color of water
pH
Taste and odor (geosmin, 2-Methylisoborneol (MIB), etc.)
Dissolved metals and salts (sodium, chloride, potassium, calcium, manganese, magnesium)
Microorganisms such as fecal coliform bacteria (Escherichia coli), Cryptosporidium, and Giardia lamblia; see Bacteriological water analysis
Dissolved metals and metalloids (lead, mercury, arsenic, etc.)
Dissolved organics: colored dissolved organic matter (CDOM), dissolved organic carbon (DOC)
Radon
Heavy metals
Pharmaceuticals
Hormone analogs
Environmental indicators
Physical indicators
Water temperature
Specific conductance or electrical conductance (EC) or conductivity
Total suspended solids (TSS)
Transparency or turbidity
Water clarity
Total dissolved solids (TDS)
Odour of water
Color of water (such as Forel-Ule scale or Pt/Co scale)
Taste of water
Chemical indicators
Alkalinity
Biochemical oxygen demand (BOD)
Chemical oxygen demand (COD)
Dissolved oxygen (DO)
Total hardness (TH)
Heavy metals
Nitrate
Orthophosphates
pH
Pesticides
Residual sodium carbonate index (RSC)
Sodium adsorption ratio (SAR)
Surfactants
Biological indicators
Ephemeroptera
Plecoptera
Mollusca
Trichoptera
Escherichia coli (E. coli)
Coliform bacteria
Pimephales promelas (fathead minnow)
Americamysis bahia (Mysid shrimp)
Sea urchin
Protists, e.g. Paratrimastix pyriformis
Biological monitoring metrics have been developed in many places, and one widely used family of measurements for freshwater is the presence and abundance of members of the insect orders Ephemeroptera, Plecoptera and Trichoptera (EPT) (of benthic macroinvertebrates whose common names are, respectively, mayfly, stonefly and caddisfly). EPT indexes will naturally vary from region to region, but generally, within a region, the greater the number of taxa from these orders, the better the water quality. Organisations in the United States, such as EPA. offer guidance on developing a monitoring program and identifying members of these and other aquatic insect orders. Many US wastewater dischargers (e.g., factories, power plants, refineries, mines, municipal sewage treatment plants) are required to conduct periodic whole effluent toxicity (WET) tests.
Individuals interested in monitoring water quality who cannot afford or manage lab scale analysis can also use biological indicators to get a general reading of water quality. One example is the IOWATER volunteer water monitoring program of Iowa, which includes an EPT indicator key.
Bivalve molluscs are largely used as bioindicators to monitor the health of aquatic environments in both fresh water and the marine environments. Their population status or structure, physiology, behaviour or the level of contamination with elements or compounds can indicate the state of contamination status of the ecosystem. They are particularly useful since they are sessile so that they are representative of the environment where they are sampled or placed. A typical project is the U.S. Mussel Watch Programme, but today they are used worldwide.
The Southern African Scoring System (SASS) method is a biological water quality monitoring system based on the presence of benthic macroinvertebrates (EPT). The SASS aquatic biomonitoring tool has been refined over the past 30 years and is now on the fifth version (SASS5) which has been specifically modified in accordance with international standards, namely the ISO/IEC 17025 protocol. The SASS5 method is used by the South African Department of Water Affairs as a standard method for River Health Assessment, which feeds the national River Health Programme and the national Rivers Database.
Climate change impacts
Standards and reports
In the setting of standards, agencies make political and technical/scientific decisions based on how the water will be used. In the case of natural water bodies, agencies also make some reasonable estimate of pristine conditions. Natural water bodies will vary in response to a region's environmental conditions, whereby water composition is influenced by the surrounding geological features, sediments, and rock types, topography, hydrology, and climate. Environmental scientists and aqueous geochemists work to interpret the parameters and environmental conditions that impact the water quality of a region, which in turn helps to identify the sources and fates of contaminants. Environmental lawyers and policymakers work to define legislation with the intention that water is maintained at an appropriate quality for its identified use.
Another general perception of water quality is that of a simple property that tells whether water is polluted or not. In fact, water quality is a complex subject, in part because water is a complex medium intrinsically tied to the ecology, geology, and anthropogenic activities of a region. Industrial and commercial activities (e.g. manufacturing, mining, construction, transport) are a major cause of water pollution as are runoff from agricultural areas, urban runoff and discharge of treated and untreated sewage.
International
The World Health Organization (WHO) published updated guidelines for drinking-water quality (GDWQ) in 2017.
The International Organization for Standardization (ISO) published regulation of water quality in the section of ICS 13.060, ranging from water sampling, drinking water, industrial class water, sewage, and examination of water for chemical, physical or biological properties. ICS 91.140.60 covers the standards of water supply systems.
National specifications for ambient water and drinking water
European Union
The water policy of the European Union is primarily codified in three directives:
Directive on Urban Waste Water Treatment (91/271/EEC) of 21 May 1991 concerning discharges of municipal and some industrial wastewaters;
The Drinking Water Directive (98/83/EC) of 3 November 1998 concerning potable water quality;
Water Framework Directive (2000/60/EC) of 23 October 2000 concerning water resources management.
India
Indian Council of Medical Research (ICMR) Standards for Drinking Water.
South Africa
Water quality guidelines for South Africa are grouped according to potential user types (e.g. domestic, industrial) in the 1996 Water Quality Guidelines. Drinking water quality is subject to the South African National Standard (SANS) 241 Drinking Water Specification.
United Kingdom
In England and Wales acceptable levels for drinking water supply are listed in the "Water Supply (Water Quality) Regulations 2000."
United States
In the United States, Water Quality Standards are defined by state agencies for various water bodies, guided by the desired uses for the water body (e.g., fish habitat, drinking water supply, recreational use). The Clean Water Act (CWA) requires each governing jurisdiction (states, territories, and covered tribal entities) to submit a set of biennial reports on the quality of water in their area. These reports are known as the 303(d) and 305(b) reports, named for their respective CWA provisions, and are submitted to, and approved by, EPA. These reports are completed by the governing jurisdiction, typically a state environmental agency. EPA recommends that each state submit a single "Integrated Report" comprising its list of impaired waters and the status of all water bodies in the state. The National Water Quality Inventory Report to Congress is a general report on water quality, providing overall information about the number of miles of streams and rivers and their aggregate condition. The CWA requires states to adopt standards for each of the possible designated uses that they assign to their waters. Should evidence suggest or document that a stream, river or lake has failed to meet the water quality criteria for one or more of its designated uses, it is placed on a list of impaired waters. Once a state has placed a water body on this list, it must develop a management plan establishing Total Maximum Daily Loads (TMDLs) for the pollutant(s) impairing the use of the water. These TMDLs establish the reductions needed to fully support the designated uses.
Drinking water standards, which are applicable to public water systems, are issued by EPA under the Safe Drinking Water Act.
| Physical sciences | Water: General | Earth science |
307065 | https://en.wikipedia.org/wiki/Tissue%20engineering | Tissue engineering | Tissue engineering is a biomedical engineering discipline that uses a combination of cells, engineering, materials methods, and suitable biochemical and physicochemical factors to restore, maintain, improve, or replace different types of biological tissues. Tissue engineering often involves the use of cells placed on tissue scaffolds in the formation of new viable tissue for a medical purpose, but is not limited to applications involving cells and tissue scaffolds. While it was once categorized as a sub-field of biomaterials, having grown in scope and importance, it can be considered as a field of its own.
While most definitions of tissue engineering cover a broad range of applications, in practice, the term is closely associated with applications that repair or replace portions of or whole tissues (i.e. organs, bone, cartilage, blood vessels, bladder, skin, muscle etc.). Often, the tissues involved require certain mechanical and structural properties for proper functioning. The term has also been applied to efforts to perform specific biochemical functions using cells within an artificially-created support system (e.g. an artificial pancreas, or a bio artificial liver). The term regenerative medicine is often used synonymously with tissue engineering, although those involved in regenerative medicine place more emphasis on the use of stem cells or progenitor cells to produce tissues.
Overview
A commonly applied definition of tissue engineering, as stated by Langer and Vacanti, is "an interdisciplinary field that applies the principles of engineering and life sciences toward the development of biological substitutes that restore, maintain, or improve [Biological tissue] function or a whole organ". In addition, Langer and Vacanti also state that there are three main types of tissue engineering: cells, tissue-inducing substances, and a cells + matrix approach (often referred to as a scaffold). Tissue engineering has also been defined as "understanding the principles of tissue growth, and applying this to produce functional replacement tissue for clinical use". A further description goes on to say that an "underlying supposition of tissue engineering is that the employment of natural biology of the system will allow for greater success in developing therapeutic strategies aimed at the replacement, repair, maintenance, or enhancement of tissue function".
Developments in the multidisciplinary field of tissue engineering have yielded a novel set of tissue replacement parts and implementation strategies. Scientific advances in biomaterials, stem cells, growth and differentiation factors, and biomimetic environments have created unique opportunities to fabricate or improve existing tissues in the laboratory from combinations of engineered extracellular matrices ("scaffolds"), cells, and biologically active molecules. Among the major challenges now facing tissue engineering is the need for more complex functionality, biomechanical stability, and vascularization in laboratory-grown tissues destined for transplantation.
Etymology
The historical origin of the term is unclear as the definition of the word has changed throughout the past few decades. The term first appeared in a 1984 publication that described the organization of an endothelium-like membrane on the surface of a long-implanted, synthetic ophthalmic prosthesis.
The first modern use of the term as recognized today was in 1985 by the researcher, physiologist and bioengineer Yuan-Cheng Fung of the Engineering Research Center. He proposed the joining of the terms tissue (in reference to the fundamental relationship between cells and organs) and engineering (in reference to the field of modification of said tissues). The term was officially adopted in 1987.
History
Ancient era (pre-17th century)
A rudimentary understanding of the inner workings of human tissues may date back further than most would expect. As early as the Neolithic period, sutures were being used to close wounds and aid in healing. Later on, societies such as ancient Egypt developed better materials for sewing up wounds such as linen sutures. Around 2500 BC in ancient India, skin grafts were developed by cutting skin from the buttock and suturing it to wound sites in the ear, nose, or lips. Ancient Egyptians often would graft skin from corpses onto living humans and even attempted to use honey as a type of antibiotic and grease as a protective barrier to prevent infection. In the 1st and 2nd centuries AD, Gallo-Romans developed wrought iron implants and dental implants could be found in ancient Mayans.
Enlightenment (17th century–19th century)
While these ancient societies had developed techniques that were way ahead of their time, they still lacked a mechanistic understanding of how the body was reacting to these procedures. This mechanistic approach came along in tandem with the development of the empirical method of science pioneered by René Descartes. Sir Isaac Newton began to describe the body as a "physiochemical machine" and postured that disease was a breakdown in the machine.
In the 17th century, Robert Hooke discovered the cell and a letter from Benedict de Spinoza brought forward the idea of the homeostasis between the dynamic processes in the body. Hydra experiments performed by Abraham Trembley in the 18th century began to delve into the regenerative capabilities of cells. During the 19th century, a better understanding of how different metals reacted with the body led to the development of better sutures and a shift towards screw and plate implants in bone fixation. Further, it was first hypothesized in the mid-1800s that cell-environment interactions and cell proliferation were vital for tissue regeneration.
Modern era (20th and 21st centuries)
As time progresses and technology advances, there is a constant need for change in the approach researchers take in their studies. Tissue engineering has continued to evolve over centuries. In the beginning people used to look at and use samples directly from human or animal cadavers. Now, tissue engineers have the ability to remake many of the tissues in the body through the use of modern techniques such as microfabrication and three-dimensional bioprinting in conjunction with native tissue cells/stem cells. These advances have allowed researchers to generate new tissues in a much more efficient manner. For example, these techniques allow for more personalization which allow for better biocompatibility, decreased immune response, cellular integration, and longevity. There is no doubt that these techniques will continue to evolve, as we have continued to see microfabrication and bioprinting evolve over the past decade.
In 1960, Wichterle and Lim were the first to publish experiments on hydrogels for biomedical applications by using them in contact lens construction. Work on the field developed slowly over the next two decades, but later found traction when hydrogels were repurposed for drug delivery. In 1984, Charles Hull developed bioprinting by converting a Hewlett-Packard inkjet printer into a device capable of depositing cells in 2-D. Three dimensional (3-D) printing is a type of additive manufacturing which has since found various applications in medical engineering, due to its high precision and efficiency. With biologist James Thompson's development of first human stem cell lines in 1998 followed by transplantation of first laboratory-grown internal organs in 1999 and creation of the first bioprinter in 2003 by the University of Missouri when they printed spheroids without the need of scaffolds, 3-D bioprinting became more conventionally used in medical field than ever before. So far, scientists have been able to print mini organoids and organs-on-chips that have rendered practical insights into the functions of a human body. Pharmaceutical companies are using these models to test drugs before moving on to animal studies. However, a fully functional and structurally similar organ has not been printed yet. A team at University of Utah has reportedly printed ears and successfully transplanted those onto children born with defects that left their ears partially developed.
Today hydrogels are considered the preferred choice of bio-inks for 3-D bioprinting since they mimic cells' natural ECM while also containing strong mechanical properties capable of sustaining 3-D structures. Furthermore, hydrogels in conjunction with 3-D bioprinting allow researchers to produce different scaffolds which can be used to form new tissues or organs. 3-D printed tissues still face many challenges such as adding vasculature. Meanwhile, 3-D printing parts of tissues definitely will improve our understanding of the human body, thus accelerating both basic and clinical research.
Examples
As defined by Langer and Vacanti, examples of tissue engineering fall into one or more of three categories: "just cells," "cells and scaffold," or "tissue-inducing factors."
In vitro meat: Edible artificial animal muscle tissue cultured in vitro.
Bioartificial liver device, "Temporary Liver", Extracorporeal Liver Assist Device (ELAD): The human hepatocyte cell line (C3A line) in a hollow fiber bioreactor can mimic the hepatic function of the liver for acute instances of liver failure. A fully capable ELAD would temporarily function as an individual's liver, thus avoiding transplantation and allowing regeneration of their own liver.
Artificial pancreas: Research involves using islet cells to regulate the body's blood sugar, particularly in cases of diabetes . Biochemical factors may be used to cause human pluripotent stem cells to differentiate (turn into) cells that function similarly to beta cells, which are in an islet cell in charge of producing insulin.
Artificial bladders: Anthony Atala (Wake Forest University) has successfully implanted artificial bladders, constructed of cultured cells seeded onto a bladder-shaped scaffold, into seven out of approximately 20 human test subjects as part of a long-term experiment.
Cartilage: lab-grown cartilage, cultured in vitro on a scaffold, was successfully used as an autologous transplant to repair patients' knees.
Scaffold-free cartilage: Cartilage generated without the use of exogenous scaffold material. In this methodology, all material in the construct is cellular produced directly by the cells.
Bioartificial heart: Doris Taylor's lab constructed a biocompatible rat heart by re-cellularising a de-cellularised rat heart. This scaffold and cells were placed in a bioreactor, where it matured to become a partially or fully transplantable organ. the work was called a "landmark". The lab first stripped the cells away from a rat heart (a process called "decellularization") and then injected rat stem cells into the decellularized rat heart.
Tissue-engineered blood vessels: Blood vessels that have been grown in a lab and can be used to repair damaged blood vessels without eliciting an immune response. Tissue engineered blood vessels have been developed by many different approaches. They could be implanted as pre-seeded cellularized blood vessels, as acellular vascular grafts made with decellularized vessels or synthetic vascular grafts.
Artificial skin constructed from human skin cells embedded in a hydrogel, such as in the case of bio-printed constructs for battlefield burn repairs.
Artificial bone marrow: Bone marrow cultured in vitro to be transplanted serves as a "just cells" approach to tissue engineering.
Tissue engineered bone: A structural matrix can be composed of metals such as titanium, polymers of varying degradation rates, or certain types of ceramics. Materials are often chosen to recruit osteoblasts to aid in reforming the bone and returning biological function. Various types of cells can be added directly into the matrix to expedite the process.
Laboratory-grown penis: Decellularized scaffolds of rabbit penises were recellularised with smooth muscle and endothelial cells. The organ was then transplanted to live rabbits and functioned comparably to the native organ, suggesting potential as treatment for genital trauma.
Oral mucosa tissue engineering uses a cells and scaffold approach to replicate the 3 dimensional structure and function of oral mucosa.
Cells as building blocks
Cells are one of the main components for the success of tissue engineering approaches. Tissue engineering uses cells as strategies for creation/replacement of new tissue. Examples include fibroblasts used for skin repair or renewal, chondrocytes used for cartilage repair (MACI–FDA approved product), and hepatocytes used in liver support systems
Cells can be used alone or with support matrices for tissue engineering applications. An adequate environment for promoting cell growth, differentiation, and integration with the existing tissue is a critical factor for cell-based building blocks. Manipulation of any of these cell processes create alternative avenues for the development of new tissue (e.g., cell reprogramming - somatic cells, vascularization).
Isolation
Techniques for cell isolation depend on the cell source. Centrifugation and apheresis are techniques used for extracting cells from biofluids (e.g., blood). Whereas digestion processes, typically using enzymes to remove the extracellular matrix (ECM), are required prior to centrifugation or apheresis techniques to extract cells from tissues/organs. Trypsin and collagenase are the most common enzymes used for tissue digestion. While trypsin is temperature dependent, collagenase is less sensitive to changes in temperature.
Cell sources
Primary cells are those directly isolated from host tissue. These cells provide an ex-vivo model of cell behavior without any genetic, epigenetic, or developmental changes; making them a closer replication of in-vivo conditions than cells derived from other methods. This constraint however, can also make studying them difficult. These are mature cells, often terminally differentiated, meaning that for many cell types proliferation is difficult or impossible. Additionally, the microenvironments these cells exist in are highly specialized, often making replication of these conditions difficult.
Secondary cells A portion of cells from a primary culture is moved to a new repository/vessel to continue being cultured. Medium from the primary culture is removed, the cells that are desired to be transferred are obtained, and then cultured in a new vessel with fresh growth medium. A secondary cell culture is useful in order to ensure that cells have both the room and nutrients that they require to grow. Secondary cultures are most notably used in any scenario in which a larger quantity of cells than can be found in the primary culture is desired. Secondary cells share the constraints of primary cells (see above) but have an added risk of contamination when transferring to a new vessel.
Genetic classifications of cells
Autologous: The donor and the recipient of the cells are the same individual. Cells are harvested, cultured or stored, and then reintroduced to the host. As a result of the host's own cells being reintroduced, an antigenic response is not elicited. The body's immune system recognizes these re-implanted cells as its own, and does not target them for attack. Autologous cell dependence on host cell health and donor site morbidity may be deterrents to their use. Adipose-derived and bone marrow-derived mesenchymal stem cells are commonly autologous in nature, and can be used in a myriad of ways, from helping repair skeletal tissue to replenishing beta cells in diabetic patients.
Allogenic: Cells are obtained from the body of a donor of the same species as the recipient. While there are some ethical constraints to the use of human cells for in vitro studies (i.e. human brain tissue chimera development), the employment of dermal fibroblasts from human foreskin demonstrates an immunologically safe and thus a viable choice for allogenic tissue engineering of the skin.
Xenogenic: These cells are derived isolated cells from alternate species from the recipient. A notable example of xenogeneic tissue utilization is cardiovascular implant construction via animal cells. Chimeric human-animal farming raises ethical concerns around the potential for improved consciousness from implanting human organs in animals.
Syngeneic or isogenic: These cells describe those borne from identical genetic code. This imparts an immunologic benefit similar to autologous cell lines (see above). Autologous cells can be considered syngenic, but the classification also extends to non-autologously derived cells such as those from an identical twin, from genetically identical (cloned) research models, or induced stem cells (iSC) as related to the donor.
Stem cells
Stem cells are undifferentiated cells with the ability to divide in culture and give rise to different forms of specialized cells. Stem cells are divided into "adult" and "embryonic" stem cells according to their source. While there is still a large ethical debate related to the use of embryonic stem cells, it is thought that another alternative source – induced pluripotent stem cellsmay be useful for the repair of diseased or damaged tissues, or may be used to grow new organs.
Totipotent cells are stem cells which can divide into further stem cells or differentiate into any cell type in the body, including extra-embryonic tissue.
Pluripotent cells are stem cells which can differentiate into any cell type in the body except extra-embryonic tissue. induced pluripotent stem cells (iPSCs) are subclass of pluripotent stem cells resembling embryonic stem cells (ESCs) that have been derived from adult differentiated cells. iPSCs are created by altering the expression of transcriptional factors in adult cells until they become like embryonic stem cells.
Multipotent stem cells can be differentiated into any cell within the same class, such as blood or bone. A common example of multipotent cells is Mesenchymal stem cells (MSCs).
Scaffolds
Scaffolds are materials that have been engineered to cause desirable cellular interactions to contribute to the formation of new functional tissues for medical purposes. Cells are often 'seeded' into these structures capable of supporting three-dimensional tissue formation. Scaffolds mimic the extracellular matrix of the native tissue, recapitulating the in vivo milieu and allowing cells to influence their own microenvironments. They usually serve at least one of the following purposes: allowing cell attachment and migration, delivering and retaining cells and biochemical factors, enabling diffusion of vital cell nutrients and expressed products, and exerting certain mechanical and biological influences to modify the behaviour of the cell phase.
In 2009, an interdisciplinary team led by the thoracic surgeon Thorsten Walles implanted the first bioartificial transplant that provides an innate vascular network for post-transplant graft supply successfully into a patient awaiting tracheal reconstruction.
To achieve the goal of tissue reconstruction, scaffolds must meet some specific requirements. High porosity and adequate pore size are necessary to facilitate cell seeding and diffusion throughout the whole structure of both cells and nutrients. Biodegradability is often an essential factor since scaffolds should preferably be absorbed by the surrounding tissues without the necessity of surgical removal. The rate at which degradation occurs has to coincide as much as possible with the rate of tissue formation: this means that while cells are fabricating their own natural matrix structure around themselves, the scaffold is able to provide structural integrity within the body and eventually it will break down leaving the newly formed tissue which will take over the mechanical load. Injectability is also important for clinical uses.
Recent research on organ printing is showing how crucial a good control of the 3D environment is to ensure reproducibility of experiments and offer better results.
Materials
Material selection is an essential aspect of producing a scaffold. The materials utilized can be natural or synthetic and can be biodegradable or non-biodegradable. Additionally, they must be biocompatible, meaning that they do not cause any adverse effects to cells. Silicone, for example, is a synthetic, non-biodegradable material commonly used as a drug delivery material, while gelatin is a biodegradable, natural material commonly used in cell-culture scaffolds
The material needed for each application is different, and dependent on the desired mechanical properties of the material. Tissue engineering of long bone defects for example, will require a rigid scaffold with a compressive strength similar to that of cortical bone (100-150 MPa), which is much higher compared to a scaffold for skin regeneration.
There are a few versatile synthetic materials used for many different scaffold applications. One of these commonly used materials is polylactic acid (PLA), a synthetic polymer. PLA – polylactic acid. This is a polyester which degrades within the human body to form lactic acid, a naturally occurring chemical which is easily removed from the body. Similar materials are polyglycolic acid (PGA) and polycaprolactone (PCL): their degradation mechanism is similar to that of PLA, but PCL degrades slower and PGA degrades faster. PLA is commonly combined with PGA to create poly-lactic-co-glycolic acid (PLGA). This is especially useful because the degradation of PLGA can be tailored by altering the weight percentages of PLA and PGA: More PLA – slower degradation, more PGA – faster degradation. This tunability, along with its biocompatibility, makes it an extremely useful material for scaffold creation.
Scaffolds may also be constructed from natural materials: in particular different derivatives of the extracellular matrix have been studied to evaluate their ability to support cell growth. Protein based materials – such as collagen, or fibrin, and polysaccharidic materials- like chitosan or glycosaminoglycans (GAGs), have all proved suitable in terms of cell compatibility. Among GAGs, hyaluronic acid, possibly in combination with cross linking agents (e.g. glutaraldehyde, water-soluble carbodiimide, etc.), is one of the possible choices as scaffold material.
Due to the covalent attachment of thiol groups to these polymers, they can crosslink via disulfide bond formation. The use of thiolated polymers (thiomers) as scaffold material for tissue engineering was initially introduced at the 4th Central European Symposium on Pharmaceutical Technology in Vienna 2001. As thiomers are biocompatible, exhibit cellular mimicking properties and efficiently support proliferation and differentiation of various cell types, they are extensively used as scaffolds for tissue engineering. Furthermore thiomers such as thiolated hyaluronic acid and thiolated chitosan were shown to exhibit wound healing properties and are subject of numerous clinical trials. Additionally, a fragment of an extracellular matrix protein, such as the RGD peptide, can be coupled to a non-bioactive material to promote cell attachment. Another form of scaffold is decellularized tissue. This is a process where chemicals are used to extracts cells from tissues, leaving just the extracellular matrix. This has the benefit of a fully formed matrix specific to the desired tissue type. However, the decellurised scaffold may present immune problems with future introduced cells.
Synthesis
A number of different methods have been described in the literature for preparing porous structures to be employed as tissue engineering scaffolds. Each of these techniques presents its own advantages, but none are free of drawbacks.
Nanofiber self-assembly
Molecular self-assembly is one of the few methods for creating biomaterials with properties similar in scale and chemistry to that of the natural in vivo extracellular matrix (ECM), a crucial step toward tissue engineering of complex tissues. Moreover, these hydrogel scaffolds have shown superiority in in vivo toxicology and biocompatibility compared to traditional macro-scaffolds and animal-derived materials.
Textile technologies
These techniques include all the approaches that have been successfully employed for the preparation of non-woven meshes of different polymers. In particular, non-woven polyglycolide structures have been tested for tissue engineering applications: such fibrous structures have been found useful to grow different types of cells. The principal drawbacks are related to the difficulties in obtaining high porosity and regular pore size.
Solvent casting and particulate leaching
Solvent casting and particulate leaching (SCPL) allows for the preparation of structures with regular porosity, but with limited thickness. First, the polymer is dissolved into a suitable organic solvent (e.g. polylactic acid could be dissolved into dichloromethane), then the solution is cast into a mold filled with porogen particles. Such porogen can be an inorganic salt like sodium chloride, crystals of saccharose, gelatin spheres or paraffin spheres. The size of the porogen particles will affect the size of the scaffold pores, while the polymer to porogen ratio is directly correlated to the amount of porosity of the final structure. After the polymer solution has been cast the solvent is allowed to fully evaporate, then the composite structure in the mold is immersed in a bath of a liquid suitable for dissolving the porogen: water in the case of sodium chloride, saccharose and gelatin or an aliphatic solvent like hexane for use with paraffin. Once the porogen has been fully dissolved, a porous structure is obtained. Other than the small thickness range that can be obtained, another drawback of SCPL lies in its use of organic solvents which must be fully removed to avoid any possible damage to the cells seeded on the scaffold.
Gas foaming
To overcome the need to use organic solvents and solid porogens, a technique using gas as a porogen has been developed. First, disc-shaped structures made of the desired polymer are prepared by means of compression molding using a heated mold. The discs are then placed in a chamber where they are exposed to high pressure CO2 for several days. The pressure inside the chamber is gradually restored to atmospheric levels. During this procedure the pores are formed by the carbon dioxide molecules that abandon the polymer, resulting in a sponge-like structure. The main problems resulting from such a technique are caused by the excessive heat used during compression molding (which prohibits the incorporation of any temperature labile material into the polymer matrix) and by the fact that the pores do not form an interconnected structure.
Emulsification freeze-drying
This technique does not require the use of a solid porogen like SCPL. First, a synthetic polymer is dissolved into a suitable solvent (e.g. polylactic acid in dichloromethane) then water is added to the polymeric solution and the two liquids are mixed in order to obtain an emulsion. Before the two phases can separate, the emulsion is cast into a mold and quickly frozen by means of immersion into liquid nitrogen. The frozen emulsion is subsequently freeze-dried to remove the dispersed water and the solvent, thus leaving a solidified, porous polymeric structure. While emulsification and freeze-drying allow for a faster preparation when compared to SCPL (since it does not require a time-consuming leaching step), it still requires the use of solvents. Moreover, pore size is relatively small and porosity is often irregular. Freeze-drying by itself is also a commonly employed technique for the fabrication of scaffolds. In particular, it is used to prepare collagen sponges: collagen is dissolved into acidic solutions of acetic acid or hydrochloric acid that are cast into a mold, frozen with liquid nitrogen and then lyophilized.
Thermally induced phase separation
Similar to the previous technique, the TIPS phase separation procedure requires the use of a solvent with a low melting point that is easy to sublime. For example, dioxane could be used to dissolve polylactic acid, then phase separation is induced through the addition of a small quantity of water: a polymer-rich and a polymer-poor phase are formed. Following cooling below the solvent melting point and some days of vacuum-drying to sublime the solvent, a porous scaffold is obtained. Liquid-liquid phase separation presents the same drawbacks of emulsification/freeze-drying.
Electrospinning
Electrospinning is a highly versatile technique that can be used to produce continuous fibers ranging in diameter from a few microns to a few nanometers. In a typical electrospinning set-up, the desired scaffold material is dissolved within a solvent and placed within a syringe. This solution is fed through a needle and a high voltage is applied to the tip and to a conductive collection surface. The buildup of electrostatic forces within the solution causes it to eject a thin fibrous stream towards the oppositely charged or grounded collection surface. During this process the solvent evaporates, leaving solid fibers leaving a highly porous network. This technique is highly tunable, with variation to solvent, voltage, working distance (distance from the needle to collection surface), flow rate of solution, solute concentration, and collection surface. This allows for precise control of fiber morphology.
On a commercial level however, due to scalability reasons, there are 40 or sometimes 96 needles involved operating at once. The bottle-necks in such set-ups are: 1) Maintaining the aforementioned variables uniformly for all of the needles and 2) formation of "beads" in single fibers that we as engineers, want to be of a uniform diameter. By modifying variables such as the distance to collector, magnitude of applied voltage, or solution flow rateresearchers can dramatically change the overall scaffold architecture.
Historically, research on electrospun fibrous scaffolds dates back to at least the late 1980s when Simon showed that electrospinning could be used to produce nano- and submicron-scale fibrous scaffolds from polymer solutions specifically intended for use as in vitro cell and tissue substrates. This early use of electrospun lattices for cell culture and tissue engineering showed that various cell types would adhere to and proliferate upon polycarbonate fibers. It was noted that as opposed to the flattened morphology typically seen in 2D culture, cells grown on the electrospun fibers exhibited a more rounded 3-dimensional morphology generally observed of tissues in vivo.
CAD/CAM technologies
Because most of the above techniques are limited when it comes to the control of porosity and pore size, computer assisted design and manufacturing techniques have been introduced to tissue engineering. First, a three-dimensional structure is designed using CAD software. The porosity can be tailored using algorithms within the software. The scaffold is then realized by using ink-jet printing of polymer powders or through Fused Deposition Modeling of a polymer melt.
A 2011 study by El-Ayoubi et al. investigated "3D-plotting technique to produce (biocompatible and biodegradable) poly-L-Lactide macroporous scaffolds with two different pore sizes" via solid free-form fabrication (SSF) with computer-aided-design (CAD), to explore therapeutic articular cartilage replacement as an "alternative to conventional tissue repair". The study found the smaller the pore size paired with mechanical stress in a bioreactor (to induce in vivo-like conditions), the higher the cell viability in potential therapeutic functionality via decreasing recovery time and increasing transplant effectiveness.
Laser-assisted bioprinting
In a 2012 study, Koch et al. focused on whether Laser-assisted BioPrinting (LaBP) can be used to build multicellular 3D patterns in natural matrix, and whether the generated constructs are functioning and forming tissue. LaBP arranges small volumes of living cell suspensions in set high-resolution patterns. The investigation was successful, the researchers foresee that "generated tissue constructs might be used for in vivo testing by implanting them into animal models" (14). As of this study, only human skin tissue has been synthesized, though researchers project that "by integrating further cell types (e.g. melanocytes, Schwann cells, hair follicle cells) into the printed cell construct, the behavior of these cells in a 3D in vitro microenvironment similar to their natural one can be analyzed", which is useful for drug discovery and toxicology studies.
Self-assembled recombinant spider silk nanomembranes
Gustafsson et al. demonstrated free‐standing, bioactive membranes of cm-sized area, but only 250 nm thin, that were formed by self‐assembly of spider silk at the interface of an aqueous solution. The membranes uniquely combine nanoscale thickness, biodegradability, ultrahigh strain and strength, permeability to proteins and promote rapid cell adherence and proliferation. They demonstrated growing a coherent layer of keratinocytes. These spider silk nanomembranes have also been used to create a static in-vitro model of a blood vessel.
Tissue engineering in situ
In situ tissue regeneration is defined as the implantation of biomaterials (alone or in combination with cells and/or biomolecules) into the tissue defect, using the surrounding microenvironment of the organism as a natural bioreactor. This approach has found application in bone regeneration, allowing the formation of cell-seeded constructs directly in the operating room.
Assembly methods
A persistent problem within tissue engineering is mass transport limitations. Engineered tissues generally lack an initial blood supply, thus making it difficult for any implanted cells to obtain sufficient oxygen and nutrients to survive, or function properly.
Self-assembly
Self-assembly methods have been shown to be promising methods for tissue engineering. Self-assembly methods have the advantage of allowing tissues to develop their own extracellular matrix, resulting in tissue that better recapitulates biochemical and biomechanical properties of native tissue. Self-assembling engineered articular cartilage was introduced by Jerry Hu and Kyriacos A. Athanasiou in 2006 and applications of the process have resulted in engineered cartilage approaching the strength of native tissue. Self-assembly is a prime technology to get cells grown in a lab to assemble into three-dimensional shapes. To break down tissues into cells, researchers first have to dissolve the extracellular matrix that normally binds them together. Once cells are isolated, they must form the complex structures that make up our natural tissues.
Liquid-based template assembly
The air-liquid surface established by Faraday waves is explored as a template to assemble biological entities for bottom-up tissue engineering. This liquid-based template can be dynamically reconfigured in a few seconds, and the assembly on the template can be achieved in a scalable and parallel manner. Assembly of microscale hydrogels, cells, neuron-seeded micro-carrier beads, cell spheroids into various symmetrical and periodic structures was demonstrated with good cell viability. Formation of 3-D neural network was achieved after 14-day tissue culture.
Additive manufacturing
It might be possible to print organs, or possibly entire organisms using additive manufacturing techniques. A recent innovative method of construction uses an ink-jet mechanism to print precise layers of cells in a matrix of thermo-reversible gel. Endothelial cells, the cells that line blood vessels, have been printed in a set of stacked rings. When incubated, these fused into a tube. This technique has been referred to as "bioprinting" within the field as it involves the printing of biological components in a structure resembling the organ of focus.
The field of three-dimensional and highly accurate models of biological systems is pioneered by multiple projects and technologies including a rapid method for creating tissues and even whole organs involve a 3-D printer that can bio-print the scaffolding and cells layer by layer into a working tissue sample or organ. The device is presented in a TED talk by Dr. Anthony Atala, M.D. the Director of the Wake Forest Institute for Regenerative Medicine, and the W.H. Boyce Professor and Chair of the Department of Urology at Wake Forest University, in which a kidney is printed on stage during the seminar and then presented to the crowd. It is anticipated that this technology will enable the production of livers in the future for transplantation and theoretically for toxicology and other biological studies as well.
In 2015 Multi-Photon Processing (MPP) was employed for in vivo experiments by engineering artificial cartilage constructs. An ex vivo histological examination showed that certain pore geometry and the pre-growing of chondrocytes (Cho) prior to implantation significantly improves the performance of the created 3-D scaffolds. The achieved biocompatibility was comparable to the commercially available collagen membranes. The successful outcome of this study supports the idea that hexagonal-pore-shaped hybrid organic-inorganic micro-structured scaffolds in combination with Cho seeding may be successfully implemented for cartilage tissue engineering.
Recently, tissue engineering has advanced with a focus on vascularization. Using Two-Photon Polymerization-based additive manufacturing, synthetic 3D microvessel networks are created from tubular hydrogel structures. These networks can perfuse tissues several cubic millimeters in size, enabling long-term viability and cell growth in vitro. This innovation marks a significant step forward in tissue engineering, facilitating the development of complex human tissue models.
Scaffolding
In 2013, using a 3-D scaffolding of Matrigel in various configurations, substantial pancreatic organoids was produced in vitro. Clusters of small numbers of cells proliferated into 40,000 cells within one week. The clusters transform into cells that make either digestive enzymes or hormones like insulin, self-organizing into branched pancreatic organoids that resemble the pancreas.
The cells are sensitive to the environment, such as gel stiffness and contact with other cells. Individual cells do not thrive; a minimum of four proximate cells was required for subsequent organoid development. Modifications to the medium composition produced either hollow spheres mainly composed of pancreatic progenitors, or complex organoids that spontaneously undergo pancreatic morphogenesis and differentiation. Maintenance and expansion of pancreatic progenitors require active Notch and FGF signaling, recapitulating in vivo niche signaling interactions.
The organoids were seen as potentially offering mini-organs for drug testing and for spare insulin-producing cells.
Aside from Matrigel 3-D scaffolds, other collagen gel systems have been developed. Collagen/hyaluronic acid scaffolds have been used for modeling the mammary gland In Vitro while co-coculturing epithelial and adipocyte cells. The HyStem kit is another 3-D platform containing ECM components and hyaluronic acid that has been used for cancer research. Additionally, hydrogel constituents can be chemically modified to assist in crosslinking and enhance their mechanical properties.
Tissue culture
In many cases, creation of functional tissues and biological structures in vitro requires extensive culturing to promote survival, growth and inducement of functionality. In general, the basic requirements of cells must be maintained in culture, which include oxygen, pH, humidity, temperature, nutrients and osmotic pressure maintenance.
Tissue engineered cultures also present additional problems in maintaining culture conditions. In standard cell culture, diffusion is often the sole means of nutrient and metabolite transport. However, as a culture becomes larger and more complex, such as the case with engineered organs and whole tissues, other mechanisms must be employed to maintain the culture, such as the creation of capillary networks within the tissue.
Another issue with tissue culture is introducing the proper factors or stimuli required to induce functionality. In many cases, simple maintenance culture is not sufficient. Growth factors, hormones, specific metabolites or nutrients, chemical and physical stimuli are sometimes required. For example, certain cells respond to changes in oxygen tension as part of their normal development, such as chondrocytes, which must adapt to low oxygen conditions or hypoxia during skeletal development. Others, such as endothelial cells, respond to shear stress from fluid flow, which is encountered in blood vessels. Mechanical stimuli, such as pressure pulses seem to be beneficial to all kind of cardiovascular tissue such as heart valves, blood vessels or pericardium.
Bioreactors
In tissue engineering, a bioreactor is a device that attempts to simulate a physiological environment in order to promote cell or tissue growth in vitro. A physiological environment can consist of many different parameters such as temperature, pressure, oxygen or carbon dioxide concentration, or osmolality of fluid environment, and it can extend to all kinds of biological, chemical or mechanical stimuli. Therefore, there are systems that may include the application of forces such as electromagnetic forces, mechanical pressures, or fluid pressures to the tissue. These systems can be two- or three-dimensional setups. Bioreactors can be used in both academic and industry applications. General-use and application-specific bioreactors are also commercially available, which may provide static chemical stimulation or a combination of chemical and mechanical stimulation.
Cell proliferation and differentiation are largely influenced by mechanical and biochemical cues in the surrounding extracellular matrix environment. Bioreactors are typically developed to replicate the specific physiological environment of the tissue being grown (e.g., flex and fluid shearing for heart tissue growth). This can allow specialized cell lines to thrive in cultures replicating their native environments, but it also makes bioreactors attractive tools for culturing stem cells. A successful stem-cell-based bioreactor is effective at expanding stem cells with uniform properties and/or promoting controlled, reproducible differentiation into selected mature cell types.
There are a variety of bioreactors designed for 3D cell cultures. There are small plastic cylindrical chambers, as well as glass chambers, with regulated internal humidity and moisture specifically engineered for the purpose of growing cells in three dimensions. The bioreactor uses bioactive synthetic materials such as polyethylene terephthalate membranes to surround the spheroid cells in an environment that maintains high levels of nutrients. They are easy to open and close, so that cell spheroids can be removed for testing, yet the chamber is able to maintain 100% humidity throughout. This humidity is important to achieve maximum cell growth and function. The bioreactor chamber is part of a larger device that rotates to ensure equal cell growth in each direction across three dimensions.
QuinXell Technologies now under Quintech Life Sciences from Singapore has developed a bioreactor known as the TisXell Biaxial Bioreactor which is specially designed for the purpose of tissue engineering. It is the first bioreactor in the world to have a spherical glass chamber with biaxial rotation; specifically to mimic the rotation of the fetus in the womb; which provides a conducive environment for the growth of tissues.
Multiple forms of mechanical stimulation have also been combined into a single bioreactor. Using gene expression analysis, one academic study found that applying a combination of cyclic strain and ultrasound stimulation to pre-osteoblast cells in a bioreactor accelerated matrix maturation and differentiation. The technology of this combined stimulation bioreactor could be used to grow bone cells more quickly and effectively in future clinical stem cell therapies.
MC2 Biotek has also developed a bioreactor known as ProtoTissue that uses gas exchange to maintain high oxygen levels within the cell chamber; improving upon previous bioreactors, since the higher oxygen levels help the cell grow and undergo normal cell respiration.
Active areas of research on bioreactors includes increasing production scale and refining the physiological environment, both of which could improve the efficiency and efficacy of bioreactors in research or clinical use. Bioreactors are currently used to study, among other things, cell and tissue level therapies, cell and tissue response to specific physiological environment changes, and development of disease and injury.
Long fiber generation
In 2013, a group from the University of Tokyo developed cell laden fibers up to a meter in length and on the order of 100 μm in size. These fibers were created using a microfluidic device that forms a double coaxial laminar flow. Each 'layer' of the microfluidic device (cells seeded in ECM, a hydrogel sheath, and finally a calcium chloride solution). The seeded cells culture within the hydrogel sheath for several days, and then the sheath is removed with viable cell fibers. Various cell types were inserted into the ECM core, including myocytes, endothelial cells, nerve cell fibers, and epithelial cell fibers. This group then showed that these fibers can be woven together to fabricate tissues or organs in a mechanism similar to textile weaving. Fibrous morphologies are advantageous in that they provide an alternative to traditional scaffold design, and many organs (such as muscle) are composed of fibrous cells.
Bioartificial organs
An artificial organ is an engineered device that can be extra corporeal or implanted to support impaired or failing organ systems. Bioartificial organs are typically created with the intent to restore critical biological functions like in the replacement of diseased hearts and lungs, or provide drastic quality of life improvements like in the use of engineered skin on burn victims. While some examples of bioartificial organs are still in the research stage of development due to the limitations involved with creating functional organs, others are currently being used in clinical settings experimentally and commercially.
Lung
Extracorporeal membrane oxygenation (ECMO) machines, otherwise known as heart and lung machines, are an adaptation of cardiopulmonary bypass techniques that provide heart and lung support. It is used primarily to support the lungs for a prolonged but still temporary timeframe (1–30 days) and allow for recovery from reversible diseases. Robert Bartlett is known as the father of ECMO and performed the first treatment of a newborn using an ECMO machine in 1975.
Skin
Tissue-engineered skin is a type of bioartificial organ that is often used to treat burns, diabetic foot ulcers, or other large wounds that cannot heal well on their own. Artificial skin can be made from autografts, allografts, and xenografts. Autografted skin comes from a patient's own skin, which allows the dermis to have a faster healing rate, and the donor site can be re-harvested a few times. Allograft skin often comes from cadaver skin and is mostly used to treat burn victims. Lastly, xenografted skin comes from animals and provides a temporary healing structure for the skin. They assist in dermal regeneration, but cannot become part of the host skin. Tissue-engineered skin is now available in commercial products. Integra, originally used to only treat burns, consists of a collagen matrix and chondroitin sulfate that can be used as a skin replacement. The chondroitin sulfate functions as a component of proteoglycans, which helps to form the extracellular matrix. Integra can be repopulated and revascularized while maintaining its dermal collagen architecture, making it a bioartificial organ Dermagraft, another commercial-made tissue-engineered skin product, is made out of living fibroblasts. These fibroblasts proliferate and produce growth factors, collagen, and ECM proteins, that help build granulation tissue.
Heart
Since the number of patients awaiting a heart transplant is continuously increasing over time, and the number of patients on the waiting list surpasses the organ availability, artificial organs used as replacement therapy for terminal heart failure would help alleviate this difficulty. Artificial hearts are usually used to bridge the heart transplantation or can be applied as replacement therapy for terminal heart malfunction. The total artificial heart (TAH), first introduced by Dr. Vladimir P. Demikhov in 1937, emerged as an ideal alternative. Since then it has been developed and improved as a mechanical pump that provides long-term circulatory support and replaces diseased or damaged heart ventricles that cannot properly pump the blood, restoring thus the pulmonary and systemic flow. Some of the current TAHs include AbioCor, an FDA-approved device that comprises two artificial ventricles and their valves, and does not require subcutaneous connections, and is indicated for patients with biventricular heart failure. In 2010 SynCardia released the portable freedom driver that allows patients to have a portable device without being confined to the hospital.
Kidney
While kidney transplants are possible, renal failure is more often treated using an artificial kidney. The first artificial kidneys and the majority of those currently in use are extracorporeal, such as with hemodialysis, which filters blood directly, or peritoneal dialysis, which filters via a fluid in the abdomen. In order to contribute to the biological functions of a kidney such as producing metabolic factors or hormones, some artificial kidneys incorporate renal cells. There has been progress in the way of making these devices smaller and more transportable, or even implantable . One challenge still to be faced in these smaller devices is countering the limited volume and therefore limited filtering capabilities.
Bioscaffolds have also been introduced to provide a framework upon which normal kidney tissue can be regenerated. These scaffolds encompass natural scaffolds (e.g., decellularized kidneys, collagen hydrogel, or silk fibroin), synthetic scaffolds (e.g., poly[lactic-co-glycolic acid] or other polymers), or a combination of two or more natural and synthetic scaffolds. These scaffolds can be implanted into the body either without cell treatment or after a period of stem cell seeding and incubation. In vitro and In vivo studies are being conducted to compare and optimize the type of scaffold and to assess whether cell seeding prior to implantation adds to the viability, regeneration and effective function of the kidneys. A recent systematic review and meta-analysis compared the results of published animal studies and identified that improved outcomes are reported with the use of hybrid (mixed) scaffolds and cell seeding; however, the meta-analysis of these results were not in agreement with the evaluation of descriptive results from the review. Therefore, further studies involving larger animals and novel scaffolds, and more transparent reproduction of previous studies are advisable.
Biomimetics
Biomimetics is a field that aims to produce materials and systems that replicate those present in nature. In the context of tissue engineering, this is a common approach used by engineers to create materials for these applications that are comparable to native tissues in terms of their structure, properties, and biocompatibility. Material properties are largely dependent on physical, structural, and chemical characteristics of that material. Subsequently, a biomimetic approach to system design will become significant in material integration, and a sufficient understanding of biological processes and interactions will be necessary. Replication of biological systems and processes may also be used in the synthesis of bio-inspired materials to achieve conditions that produce the desired biological material. Therefore, if a material is synthesized having the same characteristics of biological tissues both structurally and chemically, then ideally the synthesized material will have similar properties. This technique has an extensive history originating from the idea of using natural phenomenon as design inspiration for solutions to human problems. Many modern advancements in technology have been inspired by nature and natural systems, including aircraft, automobiles, architecture, and even industrial systems. Advancements in nanotechnology initiated the application of this technique to micro- and nano-scale problems, including tissue engineering. This technique has been used to develop synthetic bone tissues, vascular technologies, scaffolding materials and integration techniques, and functionalized nanoparticles.
Constructing neural networks in soft material
In 2018, scientists at Brandeis University reported their research on soft material embedded with chemical networks which can mimic the smooth and coordinated behavior of neural tissue. This research was funded by the U.S. Army Research Laboratory. The researchers presented an experimental system of neural networks, theoretically modeled as reaction-diffusion systems. Within the networks was an array of patterned reactors, each performing the Belousov-Zhabotinsky (BZ) reaction. These reactors could function on a nanoliter scale.
The researchers state that the inspiration for their project was the movement of the blue ribbon eel. The eel's movements are controlled by electrical impulses determined by a class of neural networks called the central pattern generator. Central Pattern Generators function within the autonomic nervous system to control bodily functions such as respiration, movement, and peristalsis.
Qualities of the reactor that were designed were the network topology, boundary conditions, initial conditions, reactor volume, coupling strength, and the synaptic polarity of the reactor (whether its behavior is inhibitory or excitatory). A BZ emulsion system with a solid elastomer polydimethylsiloxane (PDMS) was designed. Both light and bromine permeable PDMS have been reported as viable methods to create a pacemaker for neural networks.
Market
The history of the tissue engineering market can be divided into three major parts. The time before the crash of the biotech market in the early 2000s, the crash and the time afterward.
Beginning
Most early progress in tissue engineering research was done in the US. This is due to less strict regulations regarding stem cell research and more available funding than in other countries. This leads to the creation of academic startups many of them coming from Harvard or MIT. Examples are BioHybrid Technologies whose founder, Bill Chick, went to Harvard Medical School and focused on the creation of artificial pancreas. Another example would be Organogenesis Inc. whose founder went to MIT and worked on skin engineering products. Other companies with links to the MIT are TEI Biosciences, Therics and Guilford Pharmaceuticals. The renewed interest in biotechnologies in the 1980s leads to many private investors investing in these new technologies even though the business models of these early startups were often not very clear and did not present a path to long term profitability. Government sponsors were more restrained in their funding as tissue engineering was considered a high-risk investment.
In the UK the market got off to a slower start even though the regulations on stem cell research were not strict as well. This is mainly due to more investors being less willing to invest in these new technologies which were considered to be high-risk investments. Another problem faced by British companies was getting the NHS to pay for their products. This especially because the NHS runs a cost-effectiveness analysis on all supported products. Novel technologies often do not do well in this respect.
In Japan, the regulatory situation was quite different. First cell cultivation was only allowed in a hospital setting and second academic scientists employed by state-owned universities were not allowed outside employment until 1998. Moreover, the Japanese authorities took longer to approve new drugs and treatments than there US and European counterparts.
For these reasons in the early days of the Japanese market, the focus was mainly on getting products that were already approved elsewhere in Japan and selling them. Contrary to the US market the early actors in Japan were mainly big firms or sub-companies of such big firms, such as J-TEC, Menicon and Terumo, and not small startups. After regulatory changes in 2014, which allowed cell cultivation outside of a hospital setting, the speed of research in Japan increased and Japanese companies also started to develop their own products.
Crash
Soon after the big boom, the first problems started to appear. There were problems getting products approved by the FDA and if they got approved there were often difficulties in getting insurance providers to pay for the products and getting it accepted by health care providers.
For example, organogenesis ran into problems marketing its product and integrating its product in the health system. This partially due to the difficulties of handling living cells and the increased difficulties faced by physicians in using these products over conventional methods.
Another example would be Advanced Tissue Sciences Dermagraft skin product which could not create a high enough demand without reimbursements from insurance providers. Reasons for this were $4000 price-tag and the circumstance that Additionally Advanced Tissue Sciences struggled to get their product known by physicians.
The above examples demonstrate how companies struggled to make profit. This, in turn, lead investors to lose patience and stopping further funding. In consequence, several Tissue Engineering companies such as Organogenesis and Advanced Tissue Sciences filed for bankruptcy in the early 2000s. At this time, these were the only ones having commercial skin products on the market.
Reemergence
The technologies of the bankrupt or struggling companies were often bought by other companies which continued the development under more conservative business models. Examples of companies who sold their products after folding were Curis and Intercytex.
Many of the companies abandoned their long-term goals of developing fully functional organs in favor of products and technologies that could turn a profit in the short run. Examples of these kinds of products are products in the cosmetic and testing industry.
In other cases such as in the case of Advanced Tissue Sciences, the founders started new companies.
In the 2010s the regulatory framework also started to facilitate faster time to market especially in the US as new centres and pathways were created by the FDA specifically aimed at products coming from living cells such as the Center for Biologics Evaluation and Research.
The first tissue engineering products started to get commercially profitable in the 2010s.
Regulation
In Europe, regulation is currently split into three areas of regulation: medical devices, medicinal products, and biologics. Tissue engineering products are often of hybrid nature, as they are often composed of cells and a supporting structure. While some products can be approved as medicinal products, others need to gain approval as medical devices. Derksen explains in her thesis that tissue engineering researchers are sometimes confronted with regulation that does not fit the characteristics of tissue engineering.
New regulatory regimes have been observed in Europe that tackle these issues. An explanation for the difficulties in finding regulatory consensus in this matter is given by a survey conducted in the UK. The authors attribute these problems to the close relatedness and overlap with other technologies such as xenotransplantation. It can therefore not be handled separately by regulatory bodies. Regulation is further complicated by the ethical controversies associated with this and related fields of research (e.g. stem cells controversy, ethics of organ transplantation). The same survey as mentioned above shows on the example of autologous cartilage transplantation that a specific technology can be regarded as 'pure' or 'polluted' by the same social actor.
Two regulatory movements are most relevant to tissue engineering in the European Union. These are Directive 2004/23/EC on standards of quality and safety for the sourcing and processing of human tissues which was adopted by the European Parliament in 2004 and a proposed Human Tissue-Engineered Products regulation. The latter was developed under the auspices of the European Commission DG Enterprise and presented in Brussels in 2004.
| Technology | Biotechnology | null |
307156 | https://en.wikipedia.org/wiki/California%20condor | California condor | The California condor (Gymnogyps californianus) is a New World vulture and the largest North American land bird. It became extinct in the wild in 1987 when all remaining wild individuals were captured, but has since been reintroduced to northern Arizona and southern Utah (including the Grand Canyon area and Zion National Park), the coastal mountains of California, and northern Baja California in Mexico. It is the only surviving member of the genus Gymnogyps, although four extinct members of the genus are also known. The species is listed by the International Union for the Conservation of Nature as Critically Endangered, and similarly considered Critically Imperiled by NatureServe.
The plumage is black with patches of white on the underside of the wings; the head is largely bald, with skin color ranging from gray on young birds to yellow and bright orange on breeding adults. Its wingspan is the widest of any North American bird, and its weight of up to nearly equals that of the trumpeter swan, the heaviest among native North American bird species. The condor is a scavenger and eats large amounts of carrion. It is one of the world's longest-living birds, with a lifespan of up to 60 years.
Condor numbers dramatically declined in the 20th century due to agricultural chemicals (DDT), poaching, lead poisoning, and habitat destruction. A conservation plan put in place by the United States government led to the capture of all the remaining wild condors by 1987, with a total population of 27 individuals. These surviving birds were bred at the San Diego Wild Animal Park and the Los Angeles Zoo. Numbers rose through captive breeding, and beginning in 1991, condors were reintroduced into the wild. Since then, their population has grown, but the California condor remains one of the world's rarest bird species. By 31 December 2023, the Fish and Wildlife Service had updated the total world population of 561. A May 2024 population estimate of 561 is provided by the non-profit Ventana Wildlife Society on their website. The condor is a significant bird to many Californian Native American groups and plays an important role in several of their traditional myths.
Taxonomy
The California condor was described by English naturalist George Shaw in 1797 as Vultur californianus; Archibald Menzies collected the type specimen "from the coast of California" during the Vancouver expedition. It was originally classified in the same genus as the Andean condor (V. gryphus), but, due to the Andean condor's slightly different markings, slightly longer wings, and tendency to kill small animals to eat, the California condor has been placed in its own monotypic genus. The generic name Gymnogyps is derived from the Greek gymnos/γυμνος "naked" or "bare", and gyps/γυψ "vulture", while the specific name californianus comes from its location in California. The word condor itself is derived from the Quechua word kuntur.
The exact taxonomic placement of the California condor and the other six species of New World vultures remains unclear. Though similar in appearance and ecological roles to Old World vultures, the New World vultures evolved from a different ancestor in a different part of the world. Just how different the two are is under debate, with some earlier authorities suggesting that the New World vultures are more closely related to storks. More recent authorities maintain their overall position in the order Falconiformes along with the Old World vultures or place them in their own order, Cathartiformes. The South American Classification Committee has removed the New World vultures from Ciconiiformes and instead placed them in Incertae sedis, but notes that a move to Falconiformes or Cathartiformes is possible.
As of the 51st Supplement (2010) of the American Ornithologists' Union, the California condor is in the family Cathartidae of the order Cathartiformes.
Evolutionary history
The genus Gymnogyps is an example of a relict distribution. During the Pleistocene Epoch, this genus was widespread across the Americas. From fossils, the Floridian Gymnogyps kofordi from the Early Pleistocene and the Peruvian Gymnogyps howardae from the Late Pleistocene have been described. A condor found in Late Pleistocene deposits on Cuba was initially described as Antillovultur varonai, but has since been recognized as another member of Gymnogyps, Gymnogyps varonai. It may even have derived from a founder population of California condors.
The California condor is the sole surviving member of Gymnogyps and has no accepted subspecies. However, there is a Late Pleistocene form that is sometimes regarded as a palaeosubspecies, Gymnogyps californianus amplus. Opinions are mixed, regarding the classification of the form as either a chronospecies or a separate species, Gymnogyps amplus. Gymnogyps amplus occurred over much of the bird's historical range – even extending into Florida – but was larger, having about the same weight as the Andean condor. This bird also had a wider bill. As the climate changed during the last ice age, the entire population became smaller until it had evolved into the Gymnogyps californianus of today, although more recent studies by Syverson question that theory.
Description
The adult California condor is a uniform black with the exception of large triangular patches or bands of white on the underside of the wings. It has gray legs and feet, an ivory-colored bill, a frill of black feathers surrounding the base of the neck, and brownish red eyes. The juvenile is mostly a mottled dark brown with blackish coloration on the head. It has mottled gray instead of white on the underside of its flight feathers.
The condor's head has little to no feathers, which helps keep it clean when feeding on carrion. The skin of the head and neck is capable of flushing noticeably in response to emotional state. The skin color varies from yellowish to a glowing reddish-orange. The birds do not have true syringeal vocalizations. They can make a few hissing or grunting sounds only heard when very close.
The female condor is smaller than the male, an exception to the rule among birds of prey (the related Andean condor is another exception). Overall length ranges from and wingspan from . Their weight ranges from , with estimations of average weight ranging from . Wingspans of up to have been reported but no wingspan over has been verified. Most measurements are from birds raised in captivity, so it is difficult to determine if major differences exist between wild and captive condors.
California condors have the largest wingspan of any North American bird. They are surpassed in both body length and weight only by the trumpeter swan and the introduced mute swan. The American white pelican and whooping crane also have longer bodies than the condor. Condors are so large that they can be mistaken for a small, distant airplane, which possibly occurs more often than that they are mistaken for other bird species.
The middle toe of the California condor's foot is greatly elongated, and the hind toe is only slightly developed. The talons of all the toes are straight and blunt and are thus more adapted to walking than gripping. This is more similar to their supposed relatives the storks than to birds of prey and Old World vultures, which use their feet as weapons or organs of prehension.
Historic range
At the time of human settlement of the Americas, the California condor was widespread across North America; condor bones from the late Pleistocene have been found at the Cutler Fossil Site in southern Florida. However, at the end of the last glacial period came the extinction of the megafauna that led to a subsequent reduction in range and population. Five hundred years ago, the California condor roamed across the American Southwest and West Coast. Faunal remains of condors have been found documented in Arizona, Nevada, New Mexico, and Texas. The Lewis and Clark Expedition of the early 19th century reported on their sighting and shooting of California condors near the mouth of the Columbia River.
In the 1970s, two Condor Observation Sites were established in the Santa Clara River Valley to host hopeful birders interested in the endangered species: one about 15 miles north of Fillmore, California, near the Sespe Wildlife Area of Los Padres National Forest, and one atop Mount Pinos, "accessible from a dirt road off the highway in from Gorman".
Habitat
The California condor lives in rocky shrubland, coniferous forest, and oak savanna. They are often found near cliffs or large trees, which they use as nesting sites. Individual birds have a huge range and have been known to travel up to in search of carrion.
There are two sanctuaries chosen because of their prime condor nesting habitat: the Sisquoc Condor Sanctuary in the San Rafael Wilderness and the Sespe Condor Sanctuary in the Los Padres National Forest.
The Los Padres Condor Range and River Protection Act of 1992 expanded existing wilderness by and designated of new wilderness that provide habitat for the condor in the Los Padres.
Ecology and behavior
The California condor's large flight muscles are not anchored by a correspondingly large sternum, which restricts them to being primarily soarers. The birds flap their wings when taking off from the ground, but after attaining a moderate elevation they largely glide, sometimes going for miles without a single flap of their wings. They have been known to fly up to speeds of and as high as . They prefer to roost on high perches from which they can launch without any major wing-flapping effort. Often, these birds are seen soaring near rock cliffs, using thermals to aid them in keeping aloft.
The California condor has a long life span, reaching up to 60 years. If it survives to adulthood, the condor has few natural threats other than humans. Because they lack a syrinx, their vocal display is limited to grunts and hisses. Condors bathe frequently and can spend hours a day preening their feathers. Condors also perform urohidrosis, or defecate on their legs, to reduce their body temperature. There is a well-developed social structure within large groups of condors, with competition to determine a pecking order decided by body language, competitive play behavior, and a variety of hisses and grunts. This social hierarchy is displayed especially when the birds feed, with the dominant birds eating before the younger ones.
Breeding
Condors begin to look for a mate when they reach sexual maturity at the age of 6. To attract a prospective mate, the male condor performs a display, in which the male turns his head red and puffs out his neck feathers. He then spreads his wings and slowly approaches the female. If the female lowers her head to accept the male, the condors become mates for life. The pair makes a simple nest in caves or on cliff clefts, especially ones with nearby roosting trees and open spaces for landing. A mated female lays one bluish-white egg every other year. Eggs are laid as early as January to as late as April. The egg weighs about and measures from in length and about in width. If the chick or egg is lost or removed, the parents "double clutch", or lay another egg to take the lost one's place. Researchers and breeders take advantage of this behavior to double the reproductive rate by taking the first egg away for puppet-rearing; this induces the parents to lay a second (or even third) egg, which the condors are sometimes allowed to raise.
The eggs hatch after 53 to 60 days of incubation by both parents. Chicks are born with their eyes open and sometimes can take up to a week to leave the shell completely. The young are covered with a grayish down until they are almost as large as their parents. They are able to fly after 5 to 6 months, but continue to roost and forage with their parents until they are in their second year, at which point the parents typically turn their energies to a new nest. Ravens are the main predatory threat to condor eggs, while golden eagles and bears are potential predators of condor offspring.
In 2021, the San Diego Zoo reported having had two unfertilized eggs hatch within its breeding program in 2001 and 2009, producing male young by parthenogenesis as indicated by genetic studies. The mothers had been housed with males and had mated before, but the offspring lacked markers of male paternity and showed all-maternal inheritance, suggesting the specific mechanism of parthenogenesis involved automixis, gametic fusion, or endomitosis. Earlier evidence of similar parthenogenesis in birds found that among the known examples the embryos died before hatching, unlike these condor chicks. Neither chick lived to sexual maturity, preventing data collection on their reproductive potential.
In July 2024, the LA Zoo reported that a record-setting 17 California condor chicks hatched during the year's breeding season, crediting the surge on novel breeding and rearing techniques developed by their condor team. The technique involves introducing 2 to 3 chicks to a single surrogate mature condor who raises them. Due to the endangered status of the California condor, all 17 chicks are to be released into the wild.
Feeding
Wild condors maintain a large home range, often traveling a day in search of carrion. It is thought that in the early days of its existence as a species, the California condor lived off the carcasses of the Pleistocene megafauna, which are largely extinct in North America. They still prefer to feast on large, terrestrial mammalian carcasses such as deer, goats, sheep, donkeys, horses, pigs, cougars, bears, or cattle. Alternatively, they may feed on the bodies of smaller mammals such as rabbits, squirrels, and coyotes, aquatic mammals such as whales and California sea lions, or salmon. Bird and reptile carcasses are rarely eaten. Condors prefer fresh kills, but they also eat decayed food when necessary. Since they do not have a sense of smell, they spot these corpses by looking for other scavengers, like eagles and smaller vultures, the latter of which cannot rip through the tougher hides of these larger animals with the efficiency of the larger condor. They can usually intimidate other scavengers away from the carcass, with the exception of bears, which will ignore them, and golden eagles, which will fight a condor over a kill or a carcass. In the wild they are intermittent eaters, often going for between a few days to two weeks without eating, then gorging themselves on of meat at once.
Conservation
The California condor conservation project may be one of the most expensive species conservation projects in United States history, costing over $35 million, including $20 million in federal and state funding, since World War II. As of 2007, the annual cost for the condor conservation program was around $2.0 million per year. Successful reintroduction of captive-bred condors into the wild has become a multi-step and complex process, fraught with the need to periodically recapture the birds to test for lead poisoning and sometimes the necessity for lead removal by chelation.
Recovery plan
As the condor's population continued to decline, discussion began about starting a captive breeding program for the birds. Opponents to this plan argued that the condors had the right to freedom and that capturing all of the condors would change the species' habits forever, and that the cost was too great. The project received the approval of the United States government, and the U.S. Fish and Wildlife service established the California Condor Recovery Program in 1979. The capture of the remaining wild condors was completed on Easter Sunday 1987, when AC-9, the last wild condor, was captured. At that point, there were only 22 surviving condors, all of them in captivity. The goal of the California Condor Recovery Plan was to establish two geographically separate populations, one in California and the other in Arizona, each with 150 birds and at least 15 breeding pairs.
The study and capture of the remaining California condors was made possible through the efforts of Jan Hamber, an ornithologist with the Santa Barbara Museum of Natural History. Hamber personally captured AC-9, the final wild California condor, and her dedication to the bird's conservation led her to compile decades of field notes into the Condor Archives, a searchable database focused on condor biology and conservation.
The captive breeding program, led by the San Diego Wild Animal Park and Los Angeles Zoo, and with other participating zoos around the country, including the Oklahoma City Zoo and Botanical Garden, got off to a slow start due to the condor's mating habits. However, utilizing the bird's ability to double clutch, biologists began removing the first egg from the nest and raising it with puppets, allowing the parents to lay another egg.
Aside from breeding programs, the Condor Recovery Center at Oakland Zoo treats condors that are ill from lead poisoning.
Reintroduction to the wild
In 1988, the United States Fish and Wildlife Service began a reintroduction experiment involving the release of captive Andean condors into the wild in California. Only females were released, to eliminate the possibility of accidentally introducing a South American species into the United States. The experiment was a success, and all the Andean condors were recaptured and re-released in South America. California condors were released in 1991 and 1992 in California at Big Sur, Pinnacles National Park and Bitter Creek National Wildlife Refuge and in 1996 at the Vermilion Cliffs release site in Arizona near the Grand Canyon. The Fish and Wildlife Service designated the Arizona condors as an experimental, nonessential animal so they would not affect land regulations or development as ranchers were concerned they could be charged with an offense if any birds were injured on their property after the release. Though the birth rate remains low in the wild, their numbers are increasing steadily through regular releases of captive-reared adolescents.
Obstacles to recovery
In modern times, numerous causes have contributed to the California condor's decline, both before and after recovery efforts began. For example, between 1992 and 2013, 237 condor deaths occurred in the wild population. The leading cause of mortality in condor nestlings is the ingestion of trash that is fed to them by their parents. Among juveniles and adults, lead poisoning (from eating animal carcasses containing lead shot) is the leading cause of death.
Significant past damage to the condor population has also been attributed to poaching, DDT poisoning, electric power lines, egg collecting, and habitat destruction. During the California Gold Rush, some condors were even kept as pets.
Reproduction
Its low clutch size (one young per nest) and late age of sexual maturity (≈6 years) make the bird vulnerable to artificial population decline.
Inbreeding may be causing increased incidence of fatal chondrodystrophic dwarfism in wild condors, as well as a syndrome presenting with 14 rather than the typical 12 tail feathers. A 2021 study found a surprising degree of genomic diversity in condors, however. Such data allow refinement to conservation strategies, helping mitigate the effects of inbreeding. One of the study's authors hopes to complete genomic analysis of all 22 individuals from which all living condors descend.
Lead poisoning
Lead poisoning is a significant threat to condors and other avian and terrestrial scavengers Fragmented lead ammunition in large game waste is highly problematic for condors due to their extremely strong digestive juices. Blood-lead analysis of wild condors showed lead isotope signature matches to ammunition purchased by researchers near the range of the affected condors. In California, the Ridley-Tree Condor Preservation Act went into effect July 1, 2008, requiring that hunters use non-lead ammunition when hunting in the condor's range. Blood lead levels in golden eagles as well as turkey vultures has declined with the implementation of the Ridley-Tree Condor Preservation Act, demonstrating that the legislation has helped reduce other species' lead exposures aside from the California condor. There is no comparable anti-lead-bullet legislation in the other states in which the condor resides.
In 2015, Bruce Rideout, director of the wildlife disease laboratories for San Diego Zoo Global, indicated that lead poisoning is the most common cause of death for juvenile and adult condors in the wild. Among wild deaths with known causes between 1992 and 2013, over 60% (excluding chicks and fledglings) have been as a result of lead poisoning. Due to condors' long lifespan (over 50 years) and relatively late age of sexual maturity (≈6 years), and small clutch size in the wild (one egg every year or two), the population is very poorly suited to withstand the neurotoxic effects of lead exposure."
According to epidemiologist Terra Kelly, until all natural food sources are free from lead-based ammunition, "lead poisoning will threaten recovery of naturally sustaining populations of condors in the wild." While researchers and veterinarians involved in the condor recovery program note that hunters who use lead-free ammunition actually provide critical sources of food for condors and other scavengers, they caution that using lead ammunition presents a serious and preventable threat to condors and other wildlife.
Other premature death
Premature condor death may also occur due to contact with golden eagles, whose talons enable defense of carrion against condors. Evidence from condor release efforts also suggests golden eagles may occasionally kill condors.
Collision with power lines can also result in condor death. Since 1994, captive-bred California condors have been trained to avoid power lines and people. Since the implementation of this aversion conditioning program, the number of condor deaths due to power lines has greatly decreased.
Trash ingestion
"Being vultures, condors not only eat dead animals but they also have been observed eating small pieces of bone [which is especially crucial during the egg-laying period]. Although extremely intelligent, condors can’t always tell the difference between small pieces of trash and pieces of bone," according to Tim Hauck, Project Director for the California Condor Reintroduction Program. Indigestible trash can cause impaction, starvation, and death if affected condors do not receive timely medical intervention. Parent birds may unintentionally feed microtrash to nestlings, which some research has shown to be the leading cause of death among wild condor nestlings.
Disease
In 2023, Highly Pathogenic Avian Influenza (HPAI) infected members of the Utah-Arizona flock, killing 21 condors (including 13 individuals from 8 breeding pairs). Other individuals were released back into the wild following medical treatment. Sixteen condors were treated as part of a vaccine trial. As of 2 February 2024, 94 condors had received at least the first of two doses of the vaccine.
During routine winter trapping intended to assess lead levels, blood samples collected from 21 condors were tested for HPAI antibodies. About half the samples showed the presence of antibodies to the H5N1 strain of HPAI, indicating these birds were exposed to the virus and survived naturally.
Population growth
Nesting milestones have been reached by the reintroduced condors. In 2003, the first nestling fledged in the wild since 1981. In March 2006, a pair of California condors, released by Ventana Wildlife Society, attempted to nest in a hollow tree near Big Sur, California. This was the first time in more than 100 years that a pair of California condors had been seen nesting in Northern California.
In October 2010, the wild condor population reached 100 individuals in its namesake state of California, plus 73 wild condors in Arizona. In November 2011, there were 394 living individuals, 205 of them in the wild and the rest in the San Diego Zoo Safari Park, the Santa Barbara Zoo, the Los Angeles Zoo, the Oregon Zoo, and the World Center for Birds of Prey in Boise, Idaho. In May 2012, the number of living individuals had reached 405, with 179 living in captivity. By June 2014, the condor population had reached 439: 225 in the wild and 214 in captivity. Official statistics from the December 2016 USFWS recorded an overall population of 446, of which 276 are wild and 170 are captive. A key milestone was reached in 2015 when more condors were born in the wild than died.
Reintroduction to Mexico
As the Recovery Program achieved milestones, a fifth active release site in Sierra de San Pedro Mártir National Park, Baja California, Mexico, was added to the three release sites in California and the release site in Arizona. In early 2007, a California condor laid an egg in Mexico for the first time since at least the 1930s.
In June 2016, three chicks that were born in Chapultepec Zoo in Mexico City, were flown to Sierra de San Pedro Mártir National Park, Baja California, Mexico. In the spring of 2009, a second wild chick was born in the Sierra de San Pedro Mártir National Park and was named Inyaa ("Sun" in the Kiliwa language) by local environmentalists.
Expanded range
In 2014, Condor #597, also known as "Lupine", was spotted near Pescadero, a coastal community south of San Francisco. Lupine had been routinely seen at Pinnacles National Park after having been released into the wild at Big Sur the previous year. Younger birds of the central California population are seeking to expand their territory, which could mean that a new range expansion is possible for the more than 60 condors flying free in central California. Also in 2014 the first successful breeding in Utah was reported. A pair of condors that had been released in Arizona, nested in Zion National Park and the hatching of one chick was confirmed. The 1,000th chick since recovery efforts began hatched in Zion in May 2019. The California condor was seen for the first time in nearly 50 years in Sequoia National Park in late May 2020.
As part of an effort headed by the Yurok tribe to reintroduce the condor (Yurok name 'prey-go-neesh') to the coastal redwoods of northern California, birds hatched at the Oregon Zoo and the World Center for Birds of Prey were released at Redwood National Park in 2022. The first condor brought to the Yurok site was called Paaytoqin from the Nez Perce language meaning 'Come back'; he is also known as 'Mentor' or #736. He was brought to the site, but not released, to help instruct the younger condors how to behave "because of his calm nature and good disposition". Mentor condors are used to serve as a role model and establish a social hierarchy within a flock as an essential part of its survival.
The first condor to be released was called Poy’-we-son (Yurok for "the one who goes ahead"), followed by Nes-kwe-chokw ("He returns"), Ney-gem’ ‘Ne-chweenkah’ ("She carries our prayers") and ‘Hlow Hoo-let’ ("At last I (or we) fly!"). The youngsters felt at home with one another having lived together at other facilities. As of March 2024 11 birds (4 females and 7 males) have been successfully introduced, with another 5 or more being released this year. An article in the North Coast Journal from November 2023 describes the 11 birds with their names and translations. By the end of November 2024, 18 condors have been released at the site.
Condor Watch
A crowdsourcing project called Condor Watch (CW) was started on April 14, 2014, and ended in 2020. Hosted by the web portal Zooniverse, volunteers were asked to examine motion-capture images of California condors associated with release sites managed by the United States Fish and Wildlife Service, National Park Service and Ventana Wildlife Society. The tasks on the website included identifying tagged condors and marking the distance to feeding sources such as animal carcasses. Biologists can then use this data to deduce which birds are at risk of lead poisoning.
Condor Watch enabled volunteers, or citizen scientists, to participate in active research. The project had up 175,000 images to view and assess, far more than the team could hope to view on their own. Lead scientist Myra Finkelstein believes volunteering is fun because it allows enthusiasts to track the "biographies" of individual condors. Citizen science has long been used in ornithology, for instance in the Audubon Society's Christmas Bird Count, which began in 1900 and the breeding bird survey which began in 1966. McCaffrey (2005) believes this approach not only directly benefits ongoing projects, but will also help train aspiring ornithologists.
Relationship with humans
Throughout its historic range, the California condor has been a popular subject of mythology and an important symbol to Native Americans. Unusually, this bird takes on different roles in the storytelling of the different tribes.
The Wiyot tribe of California say that the condor recreated mankind after Above Old Man wiped humanity out with a flood. However, other tribes, such as California's Mono, view the condor as a destroyer, not a creator; they say that Condor seized humans, cut off their heads, and drained their blood so that it would flood Ground Squirrel's home. Condor then seized Ground Squirrel after he fled, but Ground Squirrel managed to cut off Condor's head when Condor paused to take a drink of the blood. According to the Yokuts people, the condor sometimes ate the moon, causing the lunar cycle, and his wings caused eclipses. The Chumash tribe of Southern California tell that the condor was once a white bird, but it turned black when it flew too close to a fire.
Condor bones have been found in Native American graves, as have condor feather headdresses. Cave paintings of condors have also been discovered. Some tribes ritually killed condors to make ceremonial clothing out of their feathers. Shamans then danced while wearing these to reach the upper and lower spiritual worlds. Whenever a Shaman died, his clothes were said to be cursed, so new clothing had to be made for his successor.
Some researchers such as Snyder believe that this practice of making ceremonial clothing contributed to the condor's decline, writing that California Indians killed up to 700 condors each year. Snyder continues that this figure of 700 is "no doubt an unrealistically high estimate", writing that any estimate "would remain impressively high even if divided by 10". A few tribes were known to have killed condors such as the Miwok, the Patwin, the Luiseño and the Pomo but how many is not known and difficult to judge. Using available information, Wilbur writes that "a pre-European loss of condors to Indians might not have exceeded a dozen [12] or so annually." Wilbur concludes that Indians might have contributed to the decline of California condors, "but their impact was minor except in highly localized situations."
| Biology and health sciences | Accipitrimorphae | Animals |
307157 | https://en.wikipedia.org/wiki/Andean%20condor | Andean condor | The Andean condor (Vultur gryphus) is a South American New World vulture and is the only member of the genus Vultur. It is found in the Andes mountains and adjacent Pacific coasts of western South America. With a maximum wingspan of and weight of , the Andean condor is one of the largest flying birds in the world, and is generally considered to be the largest bird of prey in the world.
It is a large black vulture with a ruff of white feathers surrounding the base of the neck and, especially in the male, large white patches on the wings. The head and neck are nearly featherless, and are a dull red color, which may flush and therefore change color in response to the bird's emotional state. In the male, there is a wattle on the neck and a large, dark red comb or caruncle on the crown of the head. The female condor is smaller than the male, an exception to the usual sexual dimorphism seen in birds of prey.
The condor is primarily a scavenger, feeding on carrion. It prefers large carcasses, such as those of deer or cattle. It reaches sexual maturity at five or six years of age and nests at elevations of up to , generally on inaccessible rock ledges. One or two eggs are usually laid. It is one of the world's longest-living birds, with a lifespan of over 70 years in some cases.
The Andean condor is a national symbol of Bolivia, Chile, Colombia, Ecuador, and Peru and plays an important role in the folklore and mythology of the Andean regions. The Andean condor is considered vulnerable by the IUCN. It is threatened by habitat loss and by secondary poisoning from lead in carcasses killed by hunters. Captive breeding programs have been instituted in several countries.
Taxonomy and systematics
The Andean condor was described by Swedish scientist Carl Linnaeus in 1758 in the tenth edition of his Systema Naturae and retains its original binomial name of Vultur gryphus. The Andean condor is sometimes called the Argentinean condor, Bolivian condor, Chilean condor, Colombian condor, Ecuadorian condor, or Peruvian condor after one of the nations to which it is native. The generic term Vultur is directly taken from the Latin vultur or voltur, which means "vulture". Its specific epithet is derived from a variant of the Greek word γρυπός (grupós, "hook-nosed"). The word condor itself is derived from the Quechua kuntur.
The exact taxonomic placement of the Andean condor and the remaining six species of New World vultures remains unclear. Although both are similar in appearance and have similar ecological roles, the New World and Old World vultures evolved from different ancestors in different parts of the world and are not closely related. Just how different the two families are is currently under debate, with some earlier authorities suggesting that the New World vultures are more closely related to storks. More recent authorities maintain their overall position in the order Accipitriformes along with the Old World vultures or place them in their own order, Cathartiformes. The South American Classification Committee has removed the New World vultures from Ciconiiformes and instead described them as incertae sedis, but notes that a move to Falconiformes or Cathartiformes is possible.
The Andean condor is the only accepted living species of its genus, Vultur. Unlike the California condor (Gymnogyps californianus), which is known from extensive fossil remains and some additional ones of congeners, the fossil record of the Andean condor recovered to date is scant. Presumed Plio-Pleistocene species of South American condors were later recognized to be not different from the present species, although one known only from a few rather small bones found in a Pliocene deposit of Tarija Department, Bolivia, may have been a smaller palaeo subspecies, V. gryphus patruus.
Description
The overall length of the Andean condor can range from . Among standard measurements, the wing chord is , the tail is and the tarsus is . Measurements are usually taken from specimens reared in captivity. The mean weight is , with the males averaging about a kilogram more at , the females a kilogram less at . Condors possess the heaviest average weight for any living flying bird or animal, ahead of trumpeter swans (Cygnus buccinator) and Dalmatian pelicans (Pelecanus crispus). However, other sources claim a mean species body mass of for the Andean condor. The Andean condor is the largest living land bird capable of flight if measured in terms of average weight and wingspan, although male bustards of the largest species (far more sexually dimorphic in size) can weigh more. The mean wingspan is around and the wings have the largest surface area of any extant bird. It has a maximum wingspan of . Among living bird species, only the great albatrosses and the two largest species of pelican exceed the Andean condor in average and maximal wingspan.
The adult plumage is all black, except for a frill of white feathers at the base of the neck and, especially in the male, large white bands on the wings, which only appear after the bird's first moult. The head and neck, kept meticulously clean, are red to blackish-red, and have few feathers. Their baldness means the skin is more exposed to the sterilizing effects of dehydration and high-altitude UV light. The crown of the head is flattened, and (in the male) is topped by a dark red comb (also called a caruncle); the skin hanging from its neck is called a wattle. Males also have yellower skin. When condors are agitated (for example, during courtship), their head and neck flush, a clear signal to animals nearby. This flush of colour is especially intense in dominant males when feeding at carcasses, and can happen in just a few seconds. Juveniles are grayish-brown, but with a blackish head and neck, and a brown ruff.
The middle toe is greatly elongated, and the hind one is only slightly developed, while the talons of all the toes are comparatively straight and blunt. The feet are thus more adapted to walking, and are of little use as weapons or organs of prehension as in birds of prey and Old World vultures. The beak is hooked, and adapted to tear rotting meat. The irises of the male are brown, while those of the female are deep red. They have no eyelashes. Unlike the case with most other birds of prey, the female is smaller.
Observation of wing color patterns, and the size and shape of the male's crest, are the best ways of identifying individual Andean condors. Sighting-resighting methods assess the size and structure of populations.
Distribution and habitat
The Andean condor is found in South America in the Andes and the Santa Marta Mountains. In the north, its range begins in Venezuela and Colombia, where it is extremely rare, then continues south along the Andes in Ecuador, Peru, and Chile, through Bolivia and western Argentina to the Tierra del Fuego. In the early 19th century, the Andean condor bred from western Venezuela to Tierra del Fuego, along the entire chain of the Andes, but its range has been greatly reduced due to human activity. Its habitat is mainly composed of open grasslands and alpine areas up to in elevation. It prefers relatively open, non-forested areas which allow it to spot carrion from the air, such as the páramo or rocky, mountainous areas in general. It occasionally ranges to lowlands in eastern Bolivia, northern Peru, and southwestern Brazil, descends to lowland desert areas in Chile and Peru, and is found over southern-beech forests in Patagonia. In southern Patagonia, meadows are important for Andean condors as this habitat is likely to have herbivores present. In this region, Andean condor distributions are therefore influenced by the locations of meadows as well as cliffs for nesting and roosting.
Ecology and behavior
The condor soars with its wings held horizontally and its primary feathers bent upwards at the tips. The lack of a large sternum to anchor its correspondingly large flight muscles physiologically identifies it as primarily being a soarer. It flaps its wings on rising from the ground, but after attaining a moderate elevation it flaps its wings very rarely, relying on thermals to stay aloft. In The Voyage of the Beagle, Charles Darwin mentioned watching condors for half an hour without once observing a flap of their wings. It prefers to roost on high places from which it can launch without major wing-flapping effort. Andean condors are often seen soaring near rock cliffs, using the heat thermals to aid them in rising in the air. Flight recorders have shown that "75% of the birds' flapping was associated with take-off", and that it "flaps its wings just 1% of the time during flight". The proportion of time for flapping is more for short flights. Flapping between two thermal glides is more than flapping between two slope glides.
Like other New World vultures, the Andean condor has the unusual habit of urohidrosis: it often empties its cloaca onto its legs and feet. A cooling effect through evaporation has been proposed as a reason for this behavior, but it does not make any sense in the cold Andean habitat of the bird. Because of this habit, their legs are often streaked with a white buildup of uric acid.
There is a well-developed social structure within large groups of condors, with competition to determine a 'pecking order' by body language, competitive play behavior, and vocalizations. Generally, mature males tend to be at the top of the pecking order, with post-dispersal immature males tending to be near the bottom.
Breeding
Sexual maturity and breeding behavior do not appear in the Andean condor until the bird is five or six years of age. It may live to 50 years or more, and it mates for life. During courtship displays, the skin of the male's neck flushes, changing from dull red to bright yellow, and inflates. He approaches the female with neck outstretched, revealing the inflated neck and the chest patch, while hissing, then extends his wings and stands erect while clicking his tongue. Other courtship rituals include hissing and clucking while hopping with wings partially spread, and dancing.
The Andean condor prefers to roost and breed at elevations of . Its nest, which consists of a few sticks placed around the eggs, is created on inaccessible ledges of rock. However, in coastal areas of Peru, where there are few cliffs, some nests are simply partially shaded crannies scraped out against boulders on slopes. It deposits one bluish-white egg, weighing about and ranging from in length. Breeding occurs about every second year, in the southern Andes around October, in the central and northern Andes it can be throughout the year. The egg hatches after 54 to 58 days of incubation by both parents. If the chick or egg is lost or removed, another egg is laid to take its place. Researchers and breeders take advantage of this behavior to double the reproductive rate by taking the first egg away for hand-rearing, causing the parents to lay a second egg, which they are generally allowed to raise. The young are covered with a grayish down until they are almost as large as their parents. They are able to fly after six months, but continue to roost and hunt with their parents until age two, when they are displaced by a new clutch.
Feeding
The Andean condor is a scavenger, feeding mainly on carrion. Wild condors inhabit large territories, often traveling more than a day in search of carrion. In inland areas, they prefer large carcasses. Naturally, they feed on the largest carcasses available, which can include llamas (Lama glama), alpacas (Lama pacos), rheas (Rhea ssp.), guanacos (Lama guanicoe), deer and armadillos. Wild individuals could acquire extra carotenoids from vegetal matter contained in carcass viscera and fresh vegetation. However, most inland condors now live largely off of domestic animals, which are now more widespread in South America, such as cattle (Bos taurus), horses (Equus caballus), donkeys (Equus asinus), mules, sheep (Ovis aries), domestic pigs (Sus domesticus), domestic goats (Capra hircus) and dogs (Canis familiaris). They also feed on the carcasses of introduced game species such as wild boar (Sus scrofa), rabbits (Oryctolagus cuniculus), foxes (Vulpes vulpes) and red deer (Cervus elaphus). For condors who live around the coast, the diet consists mainly of beached carcasses of marine mammals, largely cetaceans. They will also raid the nests of smaller birds to feed on the eggs. Andean condors have been observed to do some hunting of small, live animals, such as rodents, birds and rabbits, which (given their lack of powerful, grasping feet or developed hunting technique) they usually kill by jabbing repeatedly with their bill.
Coastal areas provide a constant food supply, and in particularly plentiful areas, some Andean condors limit their foraging area to several kilometers of beach-front land. They locate carrion by spotting it or by following other scavengers, such as corvids or other vultures. It may follow New World vultures of the genus Cathartes—the turkey vulture (C. aura), the lesser yellow-headed vulture (C. burrovianus), and the greater yellow-headed vulture (C. melambrotus)—to carcasses. The Cathartes vultures forage by smell, detecting the scent of ethyl mercaptan, a gas produced by the beginnings of decay in dead animals. These smaller vultures cannot rip through the tougher hides of these larger animals with the efficiency of the larger condor, and their interactions are often an example of mutual dependence between species. However, studies have indicated that Andean condors are fairly proficient at searching out carrion without needing to rely on other scavengers to guide them to it. Black vultures (Coragyps atratus) and several mammalian carnivorous scavengers such as foxes may sometimes track Cathartes vultures for carcasses or compete with condors over available carrion but the condor is invariably dominant among the scavengers in its range. A study in Patagonia found surprisingly that condors were driving the ecology of puma (Puma concolor) in the area, apparently by routinely commandeering the powerful cat's kills (often the day following the puma's nighttime kills). It is projected that the condors were able to engage in harassment of the pumas despite the large cat's size and power, and has apparently driven the pumas to increase their kill rate in order to accommodate for their frequent losses to the scavengers. Andean condors are intermittent eaters in the wild, often going for a few days without eating, then gorging themselves on several pounds at once, sometimes to the point of being unable to lift off the ground. Because its feet and talons are not adapted to grasping, it must feed while on the ground. Like other carrion-feeders, it plays an important role in its ecosystem by disposing of carrion which would otherwise be a breeding ground for disease.
Andean condors can efficiently absorb a wide variety of carotenoid pigments from the vegetal matter within the viscera that they consume from carcasses. These include carotenoids such as β-carotene and echineone. The pigments result in the yellow skin colouration of adult males and their ability to flush their skin a brilliant yellow during contests for dominance, as well as the colour of the iris and bright orange tongues of both sexes. Captive Andean condors have a lower concentration of carotenoid pigments in their bodies than wild condors, likely because the diet of captive condors is usually restricted to just flesh. An analysis of the droppings of wild condors found that 90% contained vegetal remains, and of those that contained vegetal remains, 35% of them were composed of primarily vegetal matter (around 80% by volume). The potential sources for the vegetal matter is posited to include the viscera of herbivore carcasses as well as fresh vegetation.
Longevity
Being a slowly-maturing bird with no known natural predators in adulthood, an Andean condor is a long-lived bird. Longevity and mortality rates are not known to have been extensively studied in the wild. Some estimations of lifespans of wild birds has exceeded 50 years. In 1983, the Guinness Book of World Records considered the longest-lived bird of any species with a confirmed lifespan was an Andean condor that died after surviving 72 years in captivity, having been captured from the wild as a juvenile of undetermined age. Several species of parrot have been reported to live for perhaps over 100 years in captivity, but these (at least in 1983) were not considered authenticated. Another early captive-held specimen of condor reportedly lived for 71 years. However, these lifespans have been exceeded by a male, nicknamed "Thaao", that was kept at Beardsley Zoo in Connecticut. Thaao was born in captivity in 1930 and died on January 26, 2010, making him 79 years of age. This is the greatest verified age ever known for a bird.
Relationship with humans
Conservation status
The Andean condor is considered vulnerable by the IUCN and the Peruvian Conservation Organization. As a result of research on its plight, its status was changed to Vulnerable from Near Threatened in 2020, and only about 10,000 individuals remain. It was first placed on the United States Endangered Species list in 1970, a status which is assigned to an animal that is in danger of extinction throughout all or a significant portion of its range. Threats to its population include loss of habitat needed for foraging, secondary poisoning from animals killed by hunters and persecution. It is threatened mainly in the northern area of its range, and is extremely rare in Venezuela and Colombia, where it has undergone considerable declines in recent years. Because it is adapted to very low mortality and has correspondingly low reproductive rates, it is extremely vulnerable to human persecution, most of which stems from the fact that it is perceived as a threat by farmers due to alleged attacks on livestock. Education programs have been implemented by conservationists to dispel this misconception. Reintroduction programs using captive-bred Andean condors, which release birds hatched in North American zoos into the wild to bolster populations, have been introduced in Argentina, Venezuela, and Colombia. The first captive-bred Andean condors were released into the wild in 1989. When raising condors, human contact is minimal; chicks are fed with glove puppets which resemble adult Andean condors in order to prevent the chicks from imprinting on humans, which would endanger them upon release as they would not be wary of humans. The condors are kept in aviaries for three months prior to release, where they acclimatize to an environment similar to that which they will be released in. Released condors are tracked by satellite in order to observe their movements and to monitor whether they are still alive.
In response to the capture of all the wild individuals of the California condor, in 1988 the US Fish and Wildlife Service began a reintroduction experiment involving the release of captive Andean condors into the wild in California. Only females were released to prevent it becoming an invasive species. The experiment was a success, and all the Andean condors were recaptured and re-released in South America before the reintroduction of the California condors took place.
In June 2014, local authorities of the Ancasmarca region rescued two Andean condors that were caged and displayed in a local market as an attraction for tourists.
Role in culture
The Andean condor is a national symbol of Argentina, Bolivia, Chile, Colombia, Ecuador, Peru and Venezuelan Andes states. It is the national bird of Bolivia, Chile, Colombia, and Ecuador. It plays an important role in the folklore and mythology of the South American Andean regions, and has been represented in Andean art from onward, and they are a part of indigenous Andean religions. In Andean mythology, the Andean condor was associated with the sun deity, and was believed to be the ruler of the upper world. The Andean condor is considered a symbol of power and health by many Andean cultures, and it was believed that the bones and organs of the Andean condor possessed medicinal powers, sometimes leading to the hunting and killing of condors to obtain its bones and organs. In some versions of Peruvian bullfighting ("Yawar Fiesta" or "Blood Festival"), a condor is tied to the back of a bull, where it pecks at the animal as bullfighters fight it. The condor generally survives and is set free.
The Andean condor is a popular figure on stamps in many countries, appearing on one for Ecuador in 1958, Argentina in 1960, Peru in 1973, Bolivia in 1985, Colombia in 1992, Chile in 1935 and 2001, and Venezuela in 2004. It has also appeared on the coins and banknotes of Colombia and Chile.
| Biology and health sciences | Accipitrimorphae | Animals |
307491 | https://en.wikipedia.org/wiki/Fishery | Fishery | Fishery can mean either the enterprise of raising or harvesting fish and other aquatic life or, more commonly, the site where such enterprise takes place (a.k.a., fishing grounds). Commercial fisheries include wild fisheries and fish farms, both in freshwater waterbodies (about 10% of all catch) and the oceans (about 90%). About 500 million people worldwide are economically dependent on fisheries. 171 million tonnes of fish were produced in 2016, but overfishing is an increasing problem, causing declines in some populations.
Because of their economic and social importance, fisheries are governed by complex fisheries management practices and legal regimes that vary widely across countries. Historically, fisheries were treated with a "first-come, first-served" approach, but recent threats from human overfishing and environmental issues have required increased regulation of fisheries to prevent conflict and increase profitable economic activity on the fishery. Modern jurisdiction over fisheries is often established by a mix of international treaties and local laws.
Declining fish populations, marine pollution, and the destruction of important coastal ecosystems have introduced increasing uncertainty in important fisheries worldwide, threatening economic security and food security in many parts of the world. These challenges are further complicated by the changes in the ocean caused by climate change, which may extend the range of some fisheries while dramatically reducing the sustainability of other fisheries.
Definitions
According to the FAO, "...a fishery is an activity leading to harvesting of fish. It may involve capture of wild fish or raising of fish through aquaculture." It is typically defined in terms of the "people involved, species or type of fish, area of water or seabed, method of fishing, class of boats, purpose of the activities or a combination of the foregoing features".
The definition often includes a combination of mammal and fish fishers in a region, the latter fishing for similar species with similar gear types. Some government and private organizations, especially those focusing on recreational fishing include in their definitions not only the fishers, but the fish and habitats upon which the fish depend.
The term fish
In biology – the term fish is most strictly used to describe any aquatic vertebrate that has gills throughout life, and can also refer to those that have limbs (if any) or appendages in the shape of fish fins. Many types of aquatic animals commonly referred to as "fish" are not fish in this strict sense; examples include shellfish, cuttlefish, starfish, crayfish and jellyfish. In the strict sense, all vertebrates are cladistically fish, although colloquially "fish" is a paraphyletic term that only refers to non-tetrapod vertebrates. In earlier times, even biologists did not make any distinction — for instance, 16th century natural historians often classified seals, whales, amphibians, crocodiles and even hippopotamuses, as well as a host of marine invertebrates, as fish.
In fisheries – the term fish is used as a collective term, and includes mollusks, crustaceans and any aquatic animals that are harvested for economic value.
True fish – The biological definition of a fish (mentioned above) is sometimes called a "true fish", the vast majority of which are teleosts. True fish are also referred to as finfish or fin fish to distinguish them from other invertebrate aquatic life harvested in fisheries or aquaculture.
Types
The fishing industry which harvests fish from fisheries can be divided into three main sectors: commercial, recreational or subsistence. They can be saltwater or freshwater, wild or farmed. About 85 percent of total marine fisheries production was finfish, mainly anchoveta (4.9 million tonnes), Alaska pollock (3.4 million tonnes) and skipjack tuna (3.1 million tonnes). Examples are the salmon fishery of Alaska, the cod fishery off the Lofoten islands, the tuna fishery of the Eastern Pacific, or the shrimp farm fisheries in China. Capture fisheries can be broadly classified as industrial scale, small-scale or artisanal, and recreational.
Close to 90% of the world's fishery catches come from oceans and seas, as opposed to inland waters. These marine catches have remained relatively stable since the mid-nineties (between 80 and 86 million tonnes). Most marine fisheries are based near the coast. This is not only because harvesting from relatively shallow waters is easier than in the open ocean, but also because fish are much more abundant near the coastal shelf, due to the abundance of nutrients available there from coastal upwelling and land runoff. However, productive wild fisheries also exist in open oceans, particularly by seamounts, and inland in lakes and rivers.
Most fisheries are wild fisheries, but farmed fisheries are increasing. Farming can occur in coastal areas, such as with oyster farms, or the aquaculture of salmon, but more typically fish farming occurs inland, in lakes, ponds, tanks and other enclosures.
There are commercial fisheries worldwide for finfish, mollusks, crustaceans and echinoderms, and by extension, aquatic plants such as kelp. However, a very small number of species support the majority of the world's fisheries. Some of these species are herring, cod, anchovy, tuna, flounder, mullet, squid, shrimp, salmon, crab, lobster, oyster and scallops. All except these last four provided a worldwide catch of well over a million tonnes in 1999, with herring and sardines together providing a harvest of over 22 million metric tons in 1999. Many other species are harvested in smaller numbers.
In 2022 small-scale fisheries contribute an estimated 40 percent of the global catch and support 90 percent of the capture fisheries workforce, with women representing 40 percent. 500 million people rely on small-scale fisheries for their livelihoods, including 53 million involved in subsistence fishing, of which 45 percent are women.
In 2022 inland fisheries produced 11.3 million tonnes, harvested mainly in Asia (63.4 percent) and Africa (29.4 percent), where they are important for food security. Lead producers were India (1.9 million tonnes), Bangladesh (1.3 million tonnes), China (1.2 million tonnes), Myanmar (0.9 million tonnes) and Indonesia (0.5 million tonnes). Inland fisheries figures are likely underestimated due to the difficulties most countries face in collecting these data.
Economic importance
Directly or indirectly, the livelihood of over 500 million people in developing countries depends on fisheries and aquaculture. Overfishing, including the taking of fish beyond sustainable levels, is reducing fish stocks and employment in many world regions. It was estimated in 2014 that global fisheries were adding US$270 billion a year to global GDP, but by full implementation of sustainable fishing, that figure could rise by as much as US$50 billion. In 2022 77% of the global workforce was in Asia, 16% in Africa and 5% in Latin America and the Caribbean.
In addition to commercial and subsistence fishing, recreational (sport) fishing is popular and economically important in many regions.
Production
Total fish production in 2016 reached an all-time high of 171 million tonnes, of which 88 percent was utilized for direct human consumption, thanks to relatively stable capture fisheries production, reduced wastage and continued aquaculture growth. This production resulted in a record-high per capita consumption of 20.3 kg in 2016. Since 1961 the annual global growth in fish consumption has been twice as high as population growth. While annual growth of aquaculture has declined in recent years, significant double-digit growth is still recorded in some countries, particularly in Africa and Asia.
FAO predicted in 2018 the following major trends for the period up to 2030:
World fish production, consumption and trade are expected to increase, but with a growth rate that will slow over time.
Despite reduced capture fisheries production in China, world capture fisheries production is projected to increase slightly through increased production in other areas if resources are properly managed. Expanding world aquaculture production, although growing more slowly than in the past, is anticipated to fill the supply–demand gap.
Prices will all increase in nominal terms while declining in real terms, although remaining high.
Food fish supply will increase in all regions, while per capita fish consumption is expected to decline in Africa, which raises concerns in terms of food security.
Trade in fish and fish products is expected to increase more slowly than in the past decade, but the share of fish production that is exported is projected to remain stable.
Management
Global goals
International attention to these issues has been captured in Sustainable Development Goal 14 "Life Below Water" which sets goals for international policy focused on preserving coastal ecosystems and supporting more sustainable economic practices for coastal communities, including in their fishery and aquaculture practices.
Law
Environmental issues
Climate change
| Technology | Buildings and infrastructure | null |
307556 | https://en.wikipedia.org/wiki/Hypocotyl | Hypocotyl | The hypocotyl (short for "hypocotyledonous stem", meaning "below seed leaf") is the stem of a germinating seedling, found below the cotyledons (seed leaves) and above the radicle (root).
Eudicots
As the plant embryo grows at germination, it sends out a shoot called a radicle that becomes the primary root, and then penetrates down into the soil. After emergence of the radicle, the hypocotyl emerges and lifts the growing tip (usually including the seed coat) above the ground, bearing the embryonic leaves (called cotyledons), and the plumule that gives rise to the first true leaves. The hypocotyl is the primary organ of extension of the young plant and develops into the stem.
Monocots
The early development of a monocot seedling like cereals and other grasses is somewhat different. A structure called the coleoptile, essentially a part of the cotyledon, protects the young stem and plumule as growth pushes them up through the soil. A mesocotyl—that part of the young plant that lies between the seed (which remains buried) and the plumule—extends the shoot up to the soil surface, where secondary roots develop from just beneath the plumule. The primary root from the radicle may then fail to develop further. The mesocotyl is considered to be partly hypocotyl and partly cotyledon (see seed).
Not all monocots develop like the grasses. The onion develops in a manner similar to the first sequence described above, the seed coat and endosperm (stored food reserve) pulled upwards as the cotyledon extends. Later, the first true leaf grows from the node between the radicle and the sheath-like cotyledon, breaking through the cotyledon to grow past it.
Storage organ
In some plants, the hypocotyl becomes enlarged as a storage organ. Examples include cyclamen, gloxinia and celeriac. In cyclamen this storage organ is called a tuber.
Hypocotyl elongation assay
One of the widely used assays in the field of photobiology is the investigation of the effect of changes in light quantity and quality on hypocotyl elongation. It is frequently used to study the growth promoting vs. growth repressing effects of application of plant hormones like ethylene. Under normal light conditions, hypocotyl growth is controlled by a process called photomorphogenesis, while shading the seedlings evokes a rapid transcriptional response which negatively regulates photomorphogenesis and results in increased rates of hypocotyl growth. This rate is highest when plants are kept in darkness mediated by a process called skotomorphogenesis, which contrasts photomorphogenesis.
| Biology and health sciences | Plant anatomy and morphology: General | Biology |
307694 | https://en.wikipedia.org/wiki/European%20shag | European shag | The European shag or common shag (Gulosus aristotelis) is a species of cormorant. It is the only member of the monotypic genus Gulosus. It breeds around the rocky coasts of western and southern Europe, southwest Asia and north Africa, mainly wintering in its breeding range except for the northernmost birds. In Britain this seabird is usually referred to as simply the shag. The scientific genus name derives from the Latin for glutton. The species name aristotelis commemorates the Greek philosopher Aristotle.
Taxonomy
The European shag was formerly classified within the genus Phalacrocorax, but a 2014 study found it to be significantly more diverged than the clade containing Phalacrocorax and Urile, but basal to the clade containing Nannopterum and Leucocarbo, and thus classified it in its own genus, Gulosus. The IOC followed this classification in 2021. Gulosus is thought to have split from the Nannopterum-Leucocarbo clade between 9.0–11.2 million years ago.
Subspecies
There are three subspecies:
G. a. aristotelis – (Linnaeus, 1761): nominate, found in northwestern Europe (Atlantic Ocean coasts)
G. a. desmarestii – (Payraudeau, 1826): found in southern Europe, southwest Asia (Mediterranean and Black Sea coasts)
G. a. riggenbachi – Hartert, 1923: found in northwest African coast
The subspecies differ slightly in bill size and the breast and leg colour of young birds. Recent evidence suggests that birds on the Atlantic coast of southwest Europe are distinct from all three, and may be an as-yet undescribed subspecies.
The name shag is also used in the Southern Hemisphere for several additional species of cormorants.
Description
This is a medium-large black bird, long and with a wingspan. It has a longish tail and a yellow throat patch. Adults have a small crest in the breeding season. It is distinguished from the great cormorant by its smaller size, lighter build, thinner bill, and, in breeding adults, by the crest and metallic green-tinged sheen on the feathers. Among those differences are that a shag is smaller and has a lighter, narrower beak, and the juvenile shag has darker underparts. The European shag's tail has 12 feathers, as do the great cormorant's 14 feathers. The green sheen on the feathers results in the alternative name green cormorant sometimes being given to the European shag.
Biology
It feeds in the sea, and, unlike the great cormorant, is rare inland. It will winter along any coast that is well-supplied with fish.
The European shag is one of the deepest divers among the cormorant family. Using depth gauges, European shags recorded diving up to deep. European shags are preponderantly benthic zone feeders, i.e. they find their prey on the sea bottom. They will eat a wide range of fish but their commonest prey is the sand eel. Shags will travel many kilometres from their roosting sites in order to feed.
In UK coastal waters, dive times are typically around 20 to 45 seconds, with a recovery time of around 15 seconds between dives; this is consistent with aerobic diving, i.e. the bird depends on the oxygen in its lungs and dissolved in its bloodstream during the dive. When they dive, they jump out of the water first to give extra impetus to the dive.
It breeds on coasts, nesting on rocky ledges or in crevices or small caves. The nests are untidy heaps of rotting seaweed or twigs cemented together by the bird's own guano. The nesting season is long, beginning in late February but some nests are not started until May or even later. Three eggs are laid. Their chicks hatch without down and so they rely totally on their parents for warmth, often for a period of two months before they can fly. Fledging may occur at any time from early June to late August, exceptionally to mid-October.
Diet
The shag is a pursuit-diving seabird that feeds predominantly in benthic habitats. Due to the relative ease with which diet samples can be collected from this species (regurgitated food or pellets) and the perceived conflict between the Phalacrocoracidae and fisheries, shag diet competition has been the subject of substantial scientific interest. Evidence collected at one colony, the Isle of May, Scotland, between 1985 and 2014, suggests that shag chick diet composition in this population has diversified in response to ocean warming. Shags also feed on fewer sandeel on windy days, presumably due to the strong effect of wind on flight in this species. The year-round diet of full-grown shags at this colony has also changed over the past 3 decades, from sandeel specialists to an increasingly diverse prey base.
Distribution
The European shag can be readily seen among the following locations during the breeding season, between late April and mid-July: Saltee Islands, Ireland; Farne Islands and Isles of Scilly, England; Isle of May, Deerness and Fowlsheugh, Scotland; Runde, Norway; Iceland; Denmark; Faroe Islands; Galicia, Northern Spain; Dalmatia and Istria, Croatia. In April 2017, eight new European shags were born in Monaco.
The largest colony of European shags is in the Cíes Islands, Spain, with 2,500 pairs (25% of the world's population).
Gallery
| Biology and health sciences | Pelecanimorphae | Animals |
307795 | https://en.wikipedia.org/wiki/Ocean%20sunfish | Ocean sunfish | The ocean sunfish (Mola mola), also known as the common mola, is one of the largest bony fish in the world. It is the type species of the genus Mola, and one of five extant species in the family Molidae. It was once misidentified as the heaviest bony fish, which was actually a different and closely related species of sunfish, Mola alexandrini. Adults typically weigh between . It is native to tropical and temperate waters around the world. It resembles a fish head without a tail, and its main body is flattened laterally. Sunfish can be as tall as they are long when their dorsal and ventral fins are extended.
Many areas of sunfish biology remain poorly understood, and various research efforts are underway, including aerial surveys of populations, satellite surveillance using pop-off satellite tags, genetic analysis of tissue samples, and collection of amateur sighting data.
Adult sunfish are vulnerable to few natural predators, but sea lions, killer whales, and sharks will consume them. Sunfish are considered a delicacy in some parts of the world, including Japan, Korea, and Taiwan. In the European Union, regulations ban the sale of fish and fishery products derived from the family Molidae. Sunfish are frequently caught in gillnets.
Naming
Its common English name, sunfish, refers to the animal's habit of sunbathing at the surface of the sea. Its common names in Dutch, Portuguese, French, Spanish, Catalan, Italian, Russian, Greek, Hungarian, Norwegian, and German (maanvis, peixe lua, Poisson lune, pez luna, peix lluna, Pesce luna, рыба-луна, φεγγαρόψαρο, holdhal, månefisk and Mondfisch, respectively) mean "moon fish", in reference to its rounded shape. In German, the fish is also known as Schwimmender Kopf, or "swimming head". In Polish, it is named samogłów, meaning "head alone" or "only head", because it has no true tail. In Swedish and Danish it is known as klumpfisk, in Dutch klompvis, in Finnish möhkäkala, all of which mean "lump fish". The Chinese translation of its academic name is , meaning "toppled wheel fish". Many of the sunfish's various names allude to its flattened shape.
Taxonomy
French polymath Guillaume Rondelet wrote about the ocean sunfish in his 1554 work de Piscibus, using the term Orthagoriscus, "sucking pig" for the likeness of its body and mouth. It was originally classified in the pufferfish family as Tetraodon mola, its epithet mola is Latin for "millstone", which the fish resembles because of its gray color, rough texture, and rounded body. It is now placed in its own genus Mola and family name Molidae as the type species with two other species: Mola tecta and M. alexandrini (previously known as Mola ramsayi). Extinct relatives of Mola mola lived in the Oligocene and Miocene epochs. However, the earliest known fossil remains of Mola mola itself were found in archaeological middens dating to the Holocene epoch.
The common name "sunfish" without qualifier is used to describe the marine family Molidae and the freshwater sunfish in the family Centrarchidae, which is unrelated to Molidae. On the other hand, the name "ocean sunfish" and "mola" refer only to the family Molidae.
Description
It shares many traits common to members in the order Tetraodontiformes including pufferfish, porcupinefish, and filefish like having a beak formed from four fused teeth; sunfish fry resemble spiky pufferfish more than they resemble adult molas.
The caudal fin of the ocean sunfish is replaced by a rounded clavus, creating the body's distinct truncated shape. The body is flattened laterally, giving it a long oval shape when seen head-on. The pectoral fins are small and fan-shaped, while the dorsal fin and the anal fin are lengthened, often making the fish as tall as it is long. Specimens up to in height have been recorded.
The mature ocean sunfish has an average length of and a fin-to-fin length of . The weight of mature specimens can range from , but even larger individuals are not unheard of. The maximum size recorded is in length, and maximum weight recorded is .
The spinal column of M. mola contains fewer vertebrae and is shorter in relation to the body than that of any other fish.
Although the sunfish descended from bony ancestors, its skeleton contains largely cartilaginous tissues, which are lighter than bone, allowing it to grow to sizes impractical for other bony fishes. Its teeth are fused into a beak-like structure, which prevents them from being able to fully close their mouths, while also having pharyngeal teeth located in the throat.
The sunfish lacks a swim bladder. Some sources indicate the internal organs contain a concentrated neurotoxin, tetrodotoxin, like the organs of other poisonous tetraodontiformes, while others dispute this claim.
Fins
In the course of its evolution, the caudal fin (tail) of the sunfish disappeared, to be replaced by a lumpy pseudotail, the clavus. This structure is formed by the convergence of the dorsal and anal fins, and is used by the fish as a rudder. The smooth-denticled clavus retains 11–14 fin rays and terminates in a number of rounded ossicles.
Ocean sunfish often swim near the surface, and their protruding dorsal fins are sometimes mistaken for those of sharks. However, the two can be distinguished by the motion of the fin. Unlike most fish, the sunfish swings its dorsal fin and anal fin in a characteristic sculling motion.
Skin
Adult sunfish range from brown to silvery-grey or white, with a variety of region-specific mottled skin patterns. Coloration is often darker on the dorsal surface, fading to a lighter shade ventrally as a form of countershading camouflage. M. mola also exhibits the ability to vary skin coloration from light to dark, especially when under attack. The skin, which contains large amounts of reticulated collagen, can be up to thick on the ventral surface, and is covered by denticles and a layer of mucus instead of scales. The skin on the clavus is smoother than that on the body, where it can be as rough as sandpaper.
More than 40 species of parasites may reside on the skin and internally, motivating the fish to seek relief in a number of ways. One of the most frequent ocean sunfish parasites is the flatworm Accacoelium contortum.
In temperate regions, drifting kelp fields harbor cleaner wrasses and other fish which remove parasites from the skin of visiting sunfish. In the tropics, M. mola solicits cleaning help from reef fishes. By basking on its side at the surface, the sunfish also allows seabirds to feed on parasites from its skin. Sunfish have been reported to breach, clearing the surface by approximately , in an apparent effort to dislodge embedded parasites.
Distribution and habitat
Ocean sunfish are native to the temperate and tropical waters of every ocean in the world. Mola genotypes appear to vary widely between the Atlantic and Pacific, but genetic differences between individuals in the Northern and Southern hemispheres are minimal.
Although early research suggested sunfish moved around mainly by drifting with ocean currents (which has resulted in the sunfish sometimes being characterized as a megaplankton), individuals have been recorded swimming in a day at a cruising speed of .
They are also capable of moving rapidly when feeding or avoiding predators, to the extent that they can vertically leap out of water. Contrary to the perception that the fish spend much of their time basking at the surface, M. mola adults actually spend a large portion of their lives actively hunting at depths greater than , occupying both the epipelagic and mesopelagic zones. Sunfish are most often found in water warmer than ; prolonged periods spent in water at temperatures of or lower can lead to disorientation and eventual death. Surface basking behavior, in which a sunfish swims on its side, presenting its largest profile to the sun, may be a method of "thermally recharging" following dives into deeper, colder water in order to feed. Sightings of the fish in colder waters outside of its usual habitat, such as those southwest of England, may be evidence of increasing marine temperatures.
Sunfish are usually found alone, but occasionally in pairs.
Feeding
The diet of the ocean sunfish was formerly thought to consist primarily of various jellyfish. However, genetic analysis reveals that sunfish are actually generalist predators that consume mostly small fish, fish larvae, squid, and crustaceans, with jellyfish and salps making up only around 15% of the diet. Occasionally they will ingest eel grass. This range of food items indicates that the sunfish feeds at many levels, from the surface to deep water, and occasionally down to the seafloor in some areas.
Life cycle
Ocean sunfish may live up to ten years in captivity, but their lifespan in a natural habitat has not yet been determined. Their growth rate remains undetermined. However, a young specimen at the Monterey Bay Aquarium increased in weight from and reached a height of nearly in 15 months.
The sheer size and thick skin of an adult of the species deters many smaller predators, but younger fish are vulnerable to predation by bluefin tuna and mahi mahi. Adults are consumed by orca, sharks and sea lions.
The mating practices of the ocean sunfish are poorly understood, but spawning areas have been suggested in the North Atlantic, South Atlantic, North Pacific, South Pacific, and Indian oceans. Females of the species can produce more eggs than any other known vertebrate, up to 300 million at a time. Sunfish eggs are released into the water and externally fertilized by sperm.
Newly hatched sunfish larvae are only long and weigh less than one gram. They develop into fry that resemble miniature pufferfish, their close relatives. Sunfish fry do not have the large pectoral fins and tail fin of their adult forms, but they have body spines uncharacteristic of adult sunfish, that disappear as they grow. Young sunfish school for protection, but this behavior is abandoned as they grow. The fry that survive can grow up to 60 million times their original weight before reaching adult proportions, arguably the most extreme size growth of any vertebrate animal.
Genome
In 2016, researchers from China National Genebank and A*STAR Singapore, including Nobel laureate Sydney Brenner, sequenced the genome of the ocean sunfish and discovered several genes which might explain its fast growth rate and large body size. As member of the order Tetraodontiformes, like fugu, the sunfish has quite a compact genome, at 730 Mb in size. Analysis from this data suggests that sunfish and pufferfishes diverged approximately 68 million years ago, which corroborates the results of other recent studies based on smaller datasets.
Human interaction
Despite their size, ocean sunfish are docile and pose no threat to human divers. Injuries from sunfish are rare, although a slight danger exists from large sunfish leaping out of the water onto boats. In 2005, a -long sunfish landed on a 4-year-old boy when the fish leaped onto the boy's family's boat off the coast of Pembrokeshire, Wales. Areas where they are commonly found are popular destinations for sport dives, and sunfish at some locations have reportedly become familiar with divers. They are more of a problem to boaters than to swimmers, as they can pose a hazard to watercraft due to their large size and weight. Collisions with sunfish are common in some parts of the world and can cause damage to the hull of a boat, or to the propellers of larger ships, as well as to the fish.
The flesh of the ocean sunfish is considered a delicacy in some regions, the largest markets being Taiwan and Japan. All parts of the sunfish are used in cuisine, from the fins to the internal organs.
Some parts are used in some areas of traditional medicine. Fishery products derived from sunfish are forbidden in the European Union according to Regulation (EC) No 853/2004 of the European Parliament and of the Council, as they contain toxins that are harmful to human health.
Sunfish are accidentally but frequently caught in drift gillnet fisheries, making up nearly 30% of the total catch of the swordfish fishery employing drift gillnets in California. The bycatch rate is even higher for the Mediterranean swordfish industry, with 71% to 90% of the total catch being sunfish.
A decrease in sunfish populations may be caused by more frequent bycatch and the increasing popularity of sunfish in human diet. The fishery bycatch and destruction of ocean sunfish are unregulated worldwide. In some areas, the fish are "finned" by fishermen who regard them as worthless bait thieves; this process, in which the fins are cut off, results in the eventual death of the fish, because it can no longer propel itself without its dorsal and anal fins.
The species is also threatened by floating litter such as plastic bags which resemble jellyfish, a common prey item. Bags can choke and suffocate a fish or fill its stomach to the extent that it starves.
In art
Patterns of this fish, known as kebuku, are seen in sarongs worn by women in Lamalera, a village in the island of Lembata, in the Lesser Sunda Islands of Indonesia.
In captivity
Sunfish are not widely held in aquarium exhibits, due to the unique and demanding requirements of their care. Some Asian aquaria display them, particularly in Japan. The Kaiyukan Aquarium in Osaka is one of few aquaria with Mola mola on display, where it is reportedly as popular an attraction as the larger whale sharks.
The Lisbon Oceanarium in Portugal has ocean sunfish showcased in the main tank, and sunfish are also on display at the Denmark Nordsøen Oceanarium.
In Kamogawa Sea World the ocean sunfish named Kukey, who started captivity in 1982, set a world record for captivity for 2,993 days, living for eight years. Kukey was at the time of delivery, but was in size at the time of death.
While the first ocean sunfish to be held in an aquarium in the United States is said to have arrived at the Monterey Bay Aquarium in August 1986, other specimens have previously been held at other locations. Marineland of the Pacific, closed since 1987 and located on the Palos Verdes Peninsula in Los Angeles County, California, held at least one ocean sunfish by 1961, and in 1964 held a specimen, said to be the largest ever captured at that time. However, another specimen was brought alive to Marineland Studios Aquarium, near St. Augustine, Florida, in 1941.
Because sunfish had not been kept in captivity on a large scale before, the staff at Monterey Bay was forced to innovate and create their own methods for capture, feeding, and parasite control. By 1998, these issues were overcome, and the aquarium was able to hold a specimen for more than a year, later releasing it after its weight increased by more than 14 times. Mola mola has since become a permanent feature of the Open Sea exhibit. Monterey Bay Aquarium's largest sunfish specimen was euthanized on February 14, 2008, after an extended period of poor health.
A major concern to curators is preventive measures taken to keep specimens in captivity from injuring themselves by rubbing against the walls of a tank, since ocean sunfish cannot easily maneuver their bodies. In a smaller tank, hanging a vinyl curtain has been used as a stopgap measure to convert a cuboid tank to a rounded shape and prevent the fish from scraping against the sides. A more effective solution is simply to provide enough room for the sunfish to swim in wide circles. The tank must also be sufficiently deep to accommodate the vertical height of the sunfish, which may reach .
Feeding captive sunfish in a tank with faster-moving, more aggressive fish can also present a challenge. Eventually, the fish can be taught to respond to a floating target to be fed, and to take food from the end of a pole or from human hands.
| Biology and health sciences | Acanthomorpha | null |
307875 | https://en.wikipedia.org/wiki/Shoulder | Shoulder | The human shoulder is made up of three bones: the clavicle (collarbone), the scapula (shoulder blade), and the humerus (upper arm bone) as well as associated muscles, ligaments and tendons.
The articulations between the bones of the shoulder make up the shoulder joints. The shoulder joint, also known as the glenohumeral joint, is the major joint of the shoulder, but can more broadly include the acromioclavicular joint.
In human anatomy, the shoulder joint comprises the part of the body where the humerus attaches to the scapula, and the head sits in the glenoid cavity. The shoulder is the group of structures in the region of the joint.
The shoulder joint is the main joint of the shoulder. It is a ball and socket joint that allows the arm to rotate in a circular fashion or to hinge out and up away from the body. The joint capsule is a soft tissue envelope that encircles the glenohumeral joint and attaches to the scapula, humerus, and head of the biceps. It is lined by a thin, smooth synovial membrane. The rotator cuff is a group of four muscles that surround the shoulder joint and contribute to the shoulder's stability. The muscles of the rotator cuff are supraspinatus, subscapularis, infraspinatus, and teres minor. The cuff adheres to the glenohumeral capsule and attaches to the humeral head.
The shoulder must be mobile enough for the wide range actions of the arms and hands, but stable enough to allow for actions such as lifting, pushing, and pulling.
Structure
The shoulder consists of a ball-and-socket joint formed by the humerus and scapula and their surrounding structures - ligaments, muscles, tendons - which support the bones and maintain the relationship of one to another. These supporting structures attach to the clavicle, humerus, and scapula, the latter providing the glenoid cavity, acromion and coracoid processes. The main joint of the shoulder is the shoulder joint (or glenohumeral joint), between the humerus and the glenoid process of the scapular. The acromioclavicular joint and sternoclavicular joint also play a role in shoulder movements. White hyaline cartilage on the ends of the bones (called articular cartilage) allows the bones to glide and move on each other, and the joint space is surrounded by a synovial membrane. Around the joint space are muscles - the rotator cuff, which directly surrounds and attaches to the shoulder joint - and other muscles that help provide stability and facilitate movement.
Two filmy sac-like structures called bursae permit smooth gliding between bone, muscle, and tendon. They cushion and protect the rotator cuff from the bony arch of the acromion.
The glenoid labrum is the second kind of cartilage in the shoulder which is distinctly different from the articular cartilage. This cartilage is more fibrous or rigid than the cartilage on the ends of the ball and socket. Also, this cartilage is also found only around the socket where it is attached.
Joint
The shoulder joint (also known as the glenohumeral joint) is the main joint of the shoulder. It is a ball and socket joint that allows the arm to rotate in a circular fashion or to hinge out and up away from the body. It is formed by the articulation between the head of the humerus and the lateral scapula (specifically-the glenoid cavity of the scapula). The "ball" of the joint is the rounded, medial anterior surface of the humerus and the "socket" is formed by the glenoid cavity, the dish-shaped portion of the lateral scapula. The shallowness of the cavity and relatively loose connections between the shoulder and the rest of the body allows the arm to have tremendous mobility, at the expense of being much easier to dislocate than most other joints in the body. There is an approximately 4-to-1 disproportion in size between the large head of the humerus and the shallow glenoid cavity.The glenoid cavity is made deeper by the addition of the fibrocartilaginous ring of the glenoid labrum.
The capsule is a soft tissue envelope that encircles the glenohumeral joint and attaches to the scapula, humerus, and head of the biceps. It is lined by a thin, smooth synovial membrane. This capsule is strengthened by the coracohumeral ligament which attaches the coracoid process of the scapula to the greater tubercle of the humerus. There are also three other ligaments attaching the lesser tubercle of the humerus to lateral scapula and are collectively called the glenohumeral ligaments.
The transverse humeral ligament, which passes from the lesser tubercle to the greater tubercle of humerus, covers the intertubercular groove, in which the long head of biceps brachii travels.
Rotator cuff
The rotator cuff is an anatomical term given to the group of four muscles and their tendons that act to stabilize the shoulder. These muscles are the supraspinatus, infraspinatus, teres minor and subscapularis and that hold the head of the humerus in the glenoid cavity during movement. The cuff adheres to the glenohumeral capsule and attaches to the head of the humerus. Together, these keep the humeral head in the glenoid cavity, preventing upward migration of the humeral head caused by the pull of the deltoid muscle at the beginning of arm elevation. The infraspinatus and the teres minor, along with the anterior fibers of the deltoid muscle, are responsible for external rotation of the arm.
The four tendons of these muscles converge to form the rotator cuff tendon. This tendon, along with the articular capsule, the coracohumeral ligament, and the glenohumeral ligament complex, blend into a confluent sheet before insertion into the humeral tuberosities. The infraspinatus and teres minor fuse near their musculotendinous junctions, while the supraspinatus and subscapularis tendons join as a sheath that surrounds the biceps tendon at the entrance of the bicipital groove.
Other muscles
Muscles from the shoulder region
In addition to the four muscles of the rotator cuff, the deltoid muscle and teres major muscles arise and exist in the shoulder region itself. The deltoid muscle covers the shoulder joint on three sides, arising from the front upper third of the clavicle, the acromion, and the spine of the scapula, and travelling to insert on the deltoid tubercle of the humerus. Contraction of each part of the deltoid assists in different movements of the shoulder - flexion (clavicular part), abduction (middle part) and extension (scapular part). The teres major attaches to the outer part of the back of the scapula, beneath the teres minor, and attaches to the upper part of the humerus. It helps with medial rotation of the humerus.
Muscles from the front
Muscles from the chest wall that contribute to the shoulder are:
Muscles from the back
Armpit
The armpit () is formed by the space between the muscles of the shoulder. The nerves and blood vessels of the arm travel through the armpit, and it possesses several sets of lymph nodes that are able to be examined. The armpit is formed by the pectoralis major and minor muscles at the front, the latissimus dorsi and teres major muscles at the back, the serratus anterior muscle on its inner surface, and the intertubercular groove of the humerus on the outer side.
Nerve supply and passage
The skin around the shoulder is supplied by C2-C4 (upper), and C7 and T2 (lower area). The brachial plexus emerges as nerve roots from the cervical vertebrae C5-T1. Branches of the plexus, in particular from C5-C6, supply the majority of the muscles of the shoulder.
Blood vessels
The subclavian artery arises from the brachiocephalic trunk on the right and directly from the aorta from the left. This becomes the axillary artery as it passes beyond the first rib. The axillary artery also supplies blood to the arm, and is one of the major sources of blood to the shoulder region. The other major sources are the transverse cervical artery and the suprascapular artery, both branches of the thyrocervical trunk which itself is a branch of the subclavian artery. The blood vessels form a network (anastamosis) behind the shoulder that helps to supply blood to the arm even when the axillary artery is compromised.
Function
The muscles and joints of the shoulder allow it to move through a remarkable range of motion, making it one of the most mobile joints in the human body. The shoulder can abduct, adduct, rotate, be raised in front of and behind the torso and move through a full 360° in the sagittal plane. This tremendous range of motion also makes the shoulder extremely unstable, far more prone to dislocation and injury than other joints
The following describes the terms used for different movements of the shoulder:
Development
Puberty
Under the influence of testosterone and growth hormone, the shoulders broaden in males during puberty.
Clinical significance
The shoulder is the most mobile and potentially unstable joint in the body. Due to this, it is often prone to problems.
Fracture
Fractures of shoulder bones can include clavicular fractures, scapular fractures, and fractures of the upper humerus.
Pain
Shoulder problems, including pain, are common and can relate to any of the structures within the shoulder. The primary cause of shoulder pain is a rotator cuff tear. The supraspinatus is most commonly involved in a rotator cuff tear, but other parts of the rotator cuff may also be involved. There are different severities of a rotator cuff tear, which range from a partial tear to a full-thickness tear. A partial tear is when the tendon is thinned, but still connected to the bone. Full-thickness tears can be separated into two classes: a full-thickness incomplete tear or a full-thickness complete tear. The incomplete tear is characterized by having only a portion of the tendon disconnected from the bone, where the complete tear has the tendon completely separated off the bone. For all forms of rotator cuff tears, depending on the severity of the injury, possible treatments include rest, an arm sling, physical therapy, steroid injections, and non-steroidal anti-inflammatory drugs, or surgery.
When this type of cartilage starts to wear out (a process called arthritis), the joint becomes painful and stiff.
Arthritis
Frozen shoulder
Impingement syndrome
Shoulder dislocation
Nerve entrapment syndrome
Imaging
Imaging of the shoulder includes ultrasound, X-ray and MRI, and is guided by the suspected diagnosis and presenting symptoms.
Conventional x-rays and ultrasonography are the primary tools used to confirm a diagnosis of injuries sustained to the rotator cuff. For extended clinical questions, imaging through Magnetic Resonance with or without intraarticular contrast agent is indicated.
Hodler et al. recommend starting scanning with conventional x-rays taken from at least two planes, since this method gives a wide first impression and even has the chance of exposing any frequent shoulder pathologies, i.e., decompensated rotator cuff tears, tendinitis calcarea, dislocations, fractures, usures, and/or osteophytes. Furthermore, x-rays are required for the planning of an optimal CT or MR image.
The conventional invasive arthrography is nowadays being replaced by the non-invasive MRI and ultrasound, and is used as an imaging reserve for patients who are contraindicated for MRI, for example pacemaker-carriers with an unclear and unsure ultrasonography.
X-ray
Projectional radiography views of the shoulder include:
AP-projection 40° posterior oblique after Grashey
The body has to be rotated about 30 to 45 degrees towards the shoulder to be imaged, and the standing or sitting patient lets the arm hang. This method reveals the joint gap and the vertical alignment towards the socket.
Transaxillary projection
The arm should be abducted 80 to 100 degrees. This method reveals:
The horizontal alignment of the humerus head in respect to the socket and the lateral clavicle in respect to the acromion
Lesions of the anterior and posterior socket border, or of the tuberculum minus
The eventual non-closure of the acromial apophysis
The coraco-humeral interval
Y-projection
The lateral contour of the shoulder should be positioned in front of the film in a way that the longitudinal axis of the scapula continues parallel to the path of the rays. This method reveals:
The horizontal centralization of the humerus head and socket
The osseous margins of the coraco-acromial arch and hence the supraspinatus outlet canal
The shape of the acromion
This projection has a low tolerance for errors and, accordingly, needs proper execution. The Y-projection can be traced back to Wijnblath’s 1933 published cavitas-en-face projection.
Ultrasound
There are several advantages of ultrasound. It is relatively cheap, does not emit any radiation, is accessible, is capable of visualizing tissue function in real time, and allows the performance of provocative maneuvers in order to replicate the patient’s pain. Those benefits have helped ultrasound become a common initial choice for assessing tendons and soft tissues. Limitations include, for example, the high degree of operator dependence and the inability to define pathologies in bones. One also has to have an extensive anatomical knowledge of the examined region and keep an open mind to normal variations and artifacts created during the scan.
Although musculoskeletal ultrasound training, like medical training in general, is a lifelong process, Kissin et al. suggests that rheumatologists who taught themselves how to manipulate ultrasound can use it just as well as international musculo-skeletal ultrasound experts to diagnose common rheumatic conditions.
After the introduction of high-frequency transducers in the mid-1980s, ultrasound has become a conventional tool for taking accurate and precise images of the shoulder to support diagnosis.
Adequate for the examination are high-resolution, high-frequency transducers with a transmission frequency of 5, 7.5, and 10 MHz. To improve the focus on structures close to the skin an additional "water start-up length" is advisable. During the examination the patient is asked to be seated, the affected arm is then adducted and the elbow is bent to 90 degrees. Slow and cautious passive lateral and/or medial rotations have the effect of being able to visualize different sections of the shoulder. In order to also demonstrate those parts which are hidden under the acromion in the neutral position, a maximum medial rotation with hyperextension behind the back is required.
To avoid the different tendon echogenicities caused by different instrument settings, Middleton compared the tendon’s echogenicity with that of the deltoid muscle, which is still lege artis.
Usually the echogenicity compared to the deltoid muscle is homogeneous intensified without dorsal echo extinction. Variability with reduced or intensified echo has also been found in healthy tendons. Bilateral comparison is very helpful when distinguishing and setting boundaries between physiological variants and a possible pathological finding. Degenerative changes at the rotator cuff often are found on both sides of the body. Consequently, unilateral differences rather point to a pathological source and bilateral changes rather to a physiological variation.
In addition, a dynamic examination can help to differentiate between an ultrasound artifact and a real pathology.
To accurately evaluate the echogenicity of an ultrasound, one has to take into account the physical laws of reflection, absorption and dispersion. It is at all times important to acknowledge that the structures in the joint of the shoulder are not aligned in the transversal, coronal or sagittal plane, and that therefore during imaging of the shoulder the transducer head has to be held perpendicularly or parallel to the structures of interest. Otherwise the appearing echogenicity may not be evaluated.
MRI
Orthopedics established the MRI early on as the tool of choice for joint- and soft tissue-imaging because of its non-invasiveness, lack of radiation exposure, multi planar slicing possibilities and the high soft tissue contrast.
MRIs can provide joint details to the treating orthopedist, helping them to diagnose and decide the next appropriate therapeutic step.
To examine the shoulder, the patient should lay down with the concerned arm is in lateral rotation. For signal detection it is recommended to use a surface-coil. To find pathologies of the rotator cuff in the basic diagnostic investigation, T2-weighted sequences with fat-suppression or STIR sequences have proven value. In general, the examination should occur in the following three main planes: axial, oblique coronal and sagittal.
Most morphological changes and injuries are sustained to the supraspinatus tendon. Traumatic rotator cuff changes are often located antero-superior, meanwhile degenerative changes more likely are supero-posterior.
Tendons are predominantly composed of dense collagen fiber bundles. Because of their extreme short T2-relaxation time they appear typically signal-weak, respectively, dark. Degenerative changes, inflammations and also partial and complete tears cause loss of the original tendon structure. Fatty deposits, mucous degeneration and hemorrhages lead to an increased intratendinal T1-image. Edema formations, inflammatory changes and ruptures increase the signals in a T2-weighted image.
MRA
While using MRI, true lesions at the rotator interval region between the parts of the supraspinatus and subscapularis are all but impossible to distinguish from normal synovium and capsule.
In 1999, Weishaupt D. et al. reached through two readers a significant better visibility of pulley lesions at the rotator interval and the expected location of the reflection pulley of the long biceps and subscapularis tendon on parasagittal (reader1/reader2 sensitivity: 86%/100%; specificity: 90%/70%) and axial (reader1/reader2 sensitivity: 86%/93%; specificity: 90%/80%) MRA images.
When examining the rotator cuff, the MRA has a couple of advantages compared to the native MRI. Through a fat suppressed T2-weighted spin echo, MRA can reproduce an extreme high fat-water-contrast, which helps to detect water-deposits with better damage diagnosis in structurally changed collagen fiber bundles.
Other animals
Tetrapod forelimbs are characterised by a high degree of mobility in the shoulder-thorax connection. Lacking a solid skeletal connection between the shoulder girdle and the vertebral column, the forelimb's attachment to the trunk is instead mainly controlled by serratus lateralis and levator scapulae. Depending on locomotor style, a bone connects the shoulder girdle to the trunk in some animals; the coracoid bone in reptiles and birds, and the clavicle in primates and bats.
In primates, the shoulder shows characteristics that differ from other mammals, including a well developed clavicle, a dorsally shifted scapula with prominent acromion and spine, and a humerus featuring a straight shaft and a spherical head.
Additional images
| Biology and health sciences | Human anatomy | Health |
308158 | https://en.wikipedia.org/wiki/Vacuum%20energy | Vacuum energy | Vacuum energy is an underlying background energy that exists in space throughout the entire universe. The vacuum energy is a special case of zero-point energy that relates to the quantum vacuum.
The effects of vacuum energy can be experimentally observed in various phenomena such as spontaneous emission, the Casimir effect, and the Lamb shift, and are thought to influence the behavior of the Universe on cosmological scales. Using the upper limit of the cosmological constant, the vacuum energy of free space has been estimated to be 10−9 joules (10−2 ergs), or ~5 GeV per cubic meter. However, in quantum electrodynamics, consistency with the principle of Lorentz covariance and with the magnitude of the Planck constant suggests a much larger value of 10113 joules per cubic meter. This huge discrepancy is known as the cosmological constant problem or, colloquially, the "vacuum catastrophe."
Origin
Quantum field theory states that all fundamental fields, such as the electromagnetic field, must be quantized at every point in space. A field in physics may be envisioned as if space were filled with interconnected vibrating balls and springs, and the strength of the field is like the displacement of a ball from its rest position. The theory requires "vibrations" in, or more accurately changes in the strength of, such a field to propagate as per the appropriate wave equation for the particular field in question. The second quantization of quantum field theory requires that each such ball–spring combination be quantized, that is, that the strength of the field be quantized at each point in space. Canonically, if the field at each point in space is a simple harmonic oscillator, its quantization places a quantum harmonic oscillator at each point. Excitations of the field correspond to the elementary particles of particle physics. Thus, according to the theory, even the vacuum has a vastly complex structure and all calculations of quantum field theory must be made in relation to this model of the vacuum.
The theory considers vacuum to implicitly have the same properties as a particle, such as spin or polarization in the case of light, energy, and so on. According to the theory, most of these properties cancel out on average leaving the vacuum empty in the literal sense of the word. One important exception, however, is the vacuum energy or the vacuum expectation value of the energy. The quantization of a simple harmonic oscillator requires the lowest possible energy, or zero-point energy of such an oscillator to be
Summing over all possible oscillators at all points in space gives an infinite quantity. To remove this infinity, one may argue that only differences in energy are physically measurable, much as the concept of potential energy has been treated in classical mechanics for centuries. This argument is the underpinning of the theory of renormalization. In all practical calculations, this is how the infinity is handled.
Vacuum energy can also be thought of in terms of virtual particles (also known as vacuum fluctuations) which are created and destroyed out of the vacuum. These particles are always created out of the vacuum in particle–antiparticle pairs, which in most cases shortly annihilate each other and disappear. However, these particles and antiparticles may interact with others before disappearing, a process which can be mapped using Feynman diagrams. Note that this method of computing vacuum energy is mathematically equivalent to having a quantum harmonic oscillator at each point and, therefore, suffers the same renormalization problems.
Additional contributions to the vacuum energy come from spontaneous symmetry breaking in quantum field theory.
Implications
Vacuum energy has a number of consequences. In 1948, Dutch physicists Hendrik B. G. Casimir and Dirk Polder predicted the existence of a tiny attractive force between closely placed metal plates due to resonances in the vacuum energy in the space between them. This is now known as the Casimir effect and has since been extensively experimentally verified. It is therefore believed that the vacuum energy is "real" in the same sense that more familiar conceptual objects such as electrons, magnetic fields, etc., are real. However, alternative explanations for the Casimir effect have since been proposed.
Other predictions are harder to verify. Vacuum fluctuations are always created as particle–antiparticle pairs. The creation of these virtual particles near the event horizon of a black hole has been hypothesized by physicist Stephen Hawking to be a mechanism for the eventual "evaporation" of black holes. If one of the pair is pulled into the black hole before this, then the other particle becomes "real" and energy/mass is essentially radiated into space from the black hole. This loss is cumulative and could result in the black hole's disappearance over time. The time required is dependent on the mass of the black hole (the equations indicate that the smaller the black hole, the more rapidly it evaporates) but could be on the order of 1060 years for large solar-mass black holes.
The vacuum energy also has important consequences for physical cosmology. General relativity predicts that energy is equivalent to mass, and therefore, if the vacuum energy is "really there", it should exert a gravitational force. Essentially, a non-zero vacuum energy is expected to contribute to the cosmological constant, which affects the expansion of the universe.
Field strength of vacuum energy
The field strength of vacuum energy is a concept proposed in a theoretical study that explores the nature of the vacuum and its relationship to gravitational interactions. The study derived a mathematical framework that uses the field strength of vacuum energy as an indicator of the bulk (spacetime) resistance to localized curvature. It illustrates the association of the field strength of vacuum energy to the curvature of the background, where this concept challenges the traditional understanding of gravity and suggests that the gravitational constant, G, may not be a universal constant, but rather a parameter dependent on the field strength of vacuum energy.
Determination of the value of G has been a topic of extensive research, with numerous experiments conducted over the years in an attempt to measure its precise value. These experiments, often employing high-precision techniques, have aimed to provide accurate measurements of G and establish a consensus on its exact value. However, the outcomes of these experiments have shown significant inconsistencies, making it difficult to reach a definitive conclusion regarding the value of G. This lack of consensus has puzzled scientists and called for alternative explanations.
To test the theoretical predictions regarding the field strength of vacuum energy, specific experimental conditions involving the position of the moon are recommended in the theoretical study. These conditions aim to achieve consistent outcomes in precision measurements of G. The ultimate goal of such experiments is to either falsify or provide confirmations to the proposed theoretical framework. The significance of exploring the field strength of vacuum energy lies in its potential to revolutionize our understanding of gravity and its interactions.
History
In 1934, Georges Lemaître used an unusual perfect-fluid equation of state to interpret the cosmological constant as due to vacuum energy. In 1948, the Casimir effect provided an experimental method for a verification of the existence of vacuum energy; in 1955, however, Evgeny Lifshitz offered a different origin for the Casimir effect. In 1957, Lee and Yang proved the concepts of broken symmetry and parity violation, for which they won the Nobel prize. In 1973, Edward Tryon proposed the zero-energy universe hypothesis: that the Universe may be a large-scale quantum-mechanical vacuum fluctuation where positive mass–energy is balanced by negative gravitational potential energy. During the 1980s, there were many attempts to relate the fields that generate the vacuum energy to specific fields that were predicted by attempts at a Grand Unified Theory and to use observations of the Universe to confirm one or another version. However, the exact nature of the particles (or fields) that generate vacuum energy, with a density such as that required by inflation theory, remains a mystery.
Vacuum energy in fiction
Arthur C. Clarke's novel The Songs of Distant Earth features a starship powered by a "quantum drive" based on aspects of this theory.
In the sci-fi television/film franchise Stargate, a Zero Point Module (ZPM) is a power source that extracts zero-point energy from a micro parallel universe.
The book Star Trek: Deep Space Nine Technical Manual describes the operating principle of the so-called quantum torpedo. In this fictional weapon, an antimatter reaction is used to create a multi-dimensional membrane in a vacuum that releases at its decomposition more energy than was needed to produce it. The missing energy is removed from the vacuum. Usually about twice as much energy is released in the explosion as would correspond to the initial antimatter matter annihilation.
In the video game Half-Life 2, the item generally known as the "gravity gun" is referred to as both the "zero point field energy manipulator" and the "zero point energy field manipulator."
| Physical sciences | Particle physics: General | Physics |
308210 | https://en.wikipedia.org/wiki/Tanzanite | Tanzanite | Tanzanite is the blue and violet variety of the mineral zoisite (a calcium aluminium hydroxyl sorosilicate), caused by small amounts of vanadium. Tanzanite belongs to the epidote mineral group. Tanzanite is only found in Simanjiro District of Manyara Region in Tanzania, in a very small mining area approximately long and wide near the Mererani Hills.
Tanzanite is noted for its remarkably strong trichroism, appearing alternately blue, violet and burgundy depending on crystal orientation. Tanzanite can also appear differently when viewed under different lighting conditions. The blues appear more evident when subjected to fluorescent light and the violet hues can be seen readily when viewed under incandescent illumination. In its rough state tanzanite is coloured a reddish brown to clear, and it requires heat treatment to remove the brownish "veil" and bring out the blue violet of the stone.
The gemstone was given the name "tanzanite" by Tiffany & Co. after Tanzania, the country in which it was discovered. The scientific name of "blue-violet zoisite" was not thought to be sufficiently consumer friendly by Tiffany's marketing department, who introduced it to the market in 1968. In 2002, the American Gem Trade Association chose tanzanite as a December birthstone, the first change to their birthstone list since 1912.
Geology
Tanzanite was formed around 585 million years ago during the mid-Ediacaran Period by massive plate tectonic activity and intense heat in the area that would later become Mount Kilimanjaro. The mineral is located in a relatively complex geological environment. Deposits are typically found in the "hinge" of isoclinal folds.
Discovery
There are many accounts of the discovery of tanzanite, but only one recognised by the government of Tanzania. In January 1967, Jumanne Mhero Ngoma (originally from Same District, Kilimanjaro) stumbled upon the sparkling blue stones at the Mererani hills in the Kiteto district of the then Arusha Region (presently Manyara Region). He was issued with a certificate of recognition three years later by the then President Julius Nyerere and a financial reward of Tsh 50,000 for his efforts. In 1984, he was issued with a certificate for scientific discovery by the Tanzania Commission for Science and Technology.
Commercial history
In July 1967, Manuel de Souza, a Goan tailor and part-time gold prospector living in Arusha, found transparent fragments of blue and blue-purple gem crystals on a ridge near Mererani, some southeast of Arusha. He assumed that the mineral was olivine (peridot) but, after soon realizing it was not, he concluded it was "dumortierite" (a blue non-gem mineral). Shortly thereafter, the stones were shown to John Saul, a Nairobi-based consulting geologist and gemstone wholesaler who was then mining aquamarine in the region around Mount Kenya. Saul, who later discovered the famous ruby deposits in the Tsavo area of Kenya, eliminated dumortierite and cordierite as possibilities, and sent samples to his father, Hyman Saul, vice president at Saks Fifth Avenue in New York. Hyman Saul took the samples to the Gemological Institute of America which correctly identified the new gem as a variety of the mineral zoisite. Correct identification was also made by mineralogists at Harvard University, the British Museum and Heidelberg University, but the first person to get the identification right was Ian McCloud, a Tanzanian government geologist based in Dodoma.
Scientifically called "blue zoisite", the gemstone was renamed as tanzanite by Henry B. Platt, a great-grandson of Louis Comfort Tiffany and a vice president of Tiffany & Co., who wanted to capitalize on the rarity and single location of the gem and thought that "blue zoisite" (which might be pronounced like "blue suicide") would not sell well. Tiffany's original campaign advertised that tanzanite could now be found in two places: "in Tanzania and at Tiffany's".
From 1967, an estimated two million carats of tanzanite were mined in Tanzania before the mines were nationalized by the Tanzanian government in 1971.
Tanzanite mining developments
In 1990, the Tanzanian government split the tanzanite mines into four sections: Blocks A, B, C and D. Blocks A and C were awarded to large operators, while Blocks B and D were reserved for the local miners. In 2005 the government renewed the lease of Block C mine to TanzaniteOne, who paid US$40 million for their lease and mining license.
In June 2003, the Tanzanian government introduced legislation banning the export of unprocessed tanzanite to India. (Like many gemstones, most tanzanite is cut in Jaipur.) The reason for the ban is to attempt to spur development of local processing facilities, thereby boosting the economy and recouping profits. This ban was phased in over a two-year period, until which time only stones over 0.5 grams were affected. In 2010, the government of Tanzania banned the export of rough stones weighing more than one gram.
TanzaniteOne Mining Ltd is owned by Richland Resources, but a 2010 law in Tanzania required them to cede 50% ownership of their mining license to the Tanzanian State Mining Company (Stamico). Production in 2011 amounted to , earning them $24 million.
Following the construction of a perimeter wall around the mines, to improve security and prevent smuggling, production rose from in 2018 to a record in 2019.
On 24 June 2020, a new record for the world's largest rough tanzanite was set after a small-scale miner, Saniniu Laizer, mined two stones of and and sold them to the Government of Tanzania through Ministry of Mining for TSh 7.74 billion (US$3.35 million) surpassing a record of stone mined by TanzaniteOne in 2005.
Total reserves of tanzanite are estimated at , according to a report published in 2018. Block C, by far the largest site, has been estimated at with a Life of Mine (LOM) expected to last until the 2040s.
Factors affecting value: grading
There is no universally accepted method of grading coloured gemstones. TanzaniteOne, a major commercial player in the tanzanite market, through its non-profit subsidiary, the Tanzanite Foundation, has introduced its own color-grading system. The new system's colour-grading scales divide tanzanite colors into a range of hues, between bluish-violet, indigo and violetish-blue.
The normal primary and secondary hues in tanzanite are blue and violet. Untreated tanzanite is a trichroic gemstone, meaning that light that enters this anisotropic crystal gets refracted on different paths, with different colour absorption on each of the three optical axes. As a result of this phenomenon, a multitude of colors have been observed in various specimens: shades of purple, violet, indigo, blue, cyan, green, yellow, orange, red and brown. After heating, tanzanite becomes dichroic. The dichroic colours range from violet through bluish-violet to indigo and violetish-blue to blue.
Clarity grading in coloured gemstones is based on the eye-clean standard, that is, a gem is considered flawless if no inclusions are visible with the unaided eye (assuming 20/20 vision). The Gemological Institute of America classifies tanzanite as a Type I gemstone, meaning it is normally eye-clean. Gems with eye-visible inclusions will be traded at deep discounts.
Heat treatment
Tanzanite forms as a brownish crystal and is trichroic, which means it shows three colours – brown, blue and violet – concurrently. Heating, either underground naturally by metamorphic processes, or artificially, removes the brown or burgundy colour component to produce a stronger violet-blue color and makes the stone "dichroic", which means it only reflects blue and violet. Rarely, gem-quality tanzanite will heat to a green primary hue, almost always accompanied by a blue or violet secondary hue. These green tanzanite have some meaningful value in the collector market, but are seldom of interest to commercial buyers.
Heat-treating in a furnace is usually carried out at between for 30 minutes. The stones should not have any cracks or bubbles, as they could shatter or the cracks/ bubble could increase in size during furnace heating.
Some stones found close to the surface in the early days of the discovery (in an area now called D block) were gem-quality blue without the need for heat treatment, probably the result of a wildfire in the area which heated the stones underground. This gave rise to the idea that "D block" stones were more desirable than tanzanite found in other areas of the small tanzanite mining area.
Since heat treatment is universal, it has no effect on price, and finished gems are assumed to be heat-treated. Gemological Institute of America states that the source of heating is gemologically undetectable, but is assumed because of its prevalence.
Tanzanite may be subjected to other forms of treatment as well. Recently, coated tanzanites were discovered and tested by the AGTA and AGL laboratories. A thin layer containing cobalt, determined by X-ray fluorescence, had been applied to improve the colour. It was noted that "coatings in particular are not considered permanent", and in the United States are required to be disclosed at the point of sale.
Pleochroism in tanzanite
Pleochroism has a physical property in which the gemstone will appear to have multiple colours based on the angle of the light hitting the stone. Tanzanite is a pleochroic gemstone. Most Tanzanite are blue when viewed from one direction but can vary from violet to red when seeing from a different angle.
The physical characters can make cutting process difficult due to the problem of selecting the perfect color. The finished colour of the gemstone will vary depending on how the table cut reflects the light.
Imitation and cobalt-coated tanzanite
, tanzanite has never been successfully synthesized in a laboratory, so all genuine tanzanite is naturally occurring. However, because of its rarity and market demand, tanzanite has been imitated in several ways. Among the materials used for this are cubic zirconia, synthetic spinel, yttrium aluminium garnet, and colored glass. A test of the stone with a dichroscope can easily distinguish these from genuine tanzanite, as only tanzanite will appear doubly refractive: the two viewing windows of the dichroscope will display different colors (one window blue, the other violet) when viewing genuine tanzanite, while the imitation stones are all singly refractive and will cause both windows to appear the same color (violet).
Synthetic forsterite (, the magnesium-rich end-member of olivine) has also been sold as tanzanite, and presents a similar appearance. It can be distinguished from tanzanite in three ways. The first is by using a refractometer: forsterite will show a refractive index of between 1.63 and 1.67, while tanzanite will show a higher index of 1.685 to 1.707. The second way is by using a Hanneman filter: through it, genuine tanzanite will appear orange-pink, while forsterite will appear green. The third way is by examining a cut stone through its crown facets and viewing the pavilion cuts at the back of the stone using a standard jeweller's loupe: forsterite will show birefringence, making the pavilion cuts appear "doubled up", while the much lower birefringence of tanzanite will not have this characteristic.
Lower grades of tanzanite are occasionally enhanced using a layer of cobalt, as cobalt imparts a deeper shade of blue. The cobalt layer does not adhere well to these stones, and tends to rub off over time, resulting in a much less intensely colored stone. Though still tanzanite, the practice of cobalt coating is considered deceptive unless well-advertised.
| Physical sciences | Silicate minerals | Earth science |
308449 | https://en.wikipedia.org/wiki/Bronchus | Bronchus | A bronchus ( ; : bronchi, ) is a passage or airway in the lower respiratory tract that conducts air into the lungs. The first or primary bronchi to branch from the trachea at the carina are the right main bronchus and the left main bronchus. These are the widest bronchi, and enter the right lung, and the left lung at each hilum. The main bronchi branch into narrower secondary bronchi or lobar bronchi, and these branch into narrower tertiary bronchi or segmental bronchi. Further divisions of the segmental bronchi are known as 4th order, 5th order, and 6th order segmental bronchi, or grouped together as subsegmental bronchi.
The bronchi, when too narrow to be supported by cartilage, are known as bronchioles. No gas exchange takes place in the bronchi.
Structure
The trachea (windpipe) divides at the carina into two main or primary bronchi, the left bronchus and the right bronchus. The carina of the trachea is located at the level of the sternal angle and the fifth thoracic vertebra (at rest).
The right main bronchus is wider, shorter, and more vertical than the left main bronchus, its mean length is 1.09 cm. It enters the root of the right lung at approximately the fifth thoracic vertebra. The right main bronchus subdivides into three secondary bronchi (also known as lobar bronchi), which deliver oxygen to the three lobes of the right lung—the superior, middle and inferior lobe. The azygos vein arches over it from behind; and the right pulmonary artery lies at first below and then in front of it. About 2 cm from its commencement it gives off a branch to the superior lobe of the right lung, which is also called the eparterial bronchus. Eparterial refers to its position above the right pulmonary artery. The right bronchus now passes below the artery, and is known as the hyparterial branch which divides into the two lobar bronchi to the middle and lower lobes.
The left main bronchus is smaller in caliber but longer than the right, being 5 cm long. It enters the root of the left lung opposite the sixth thoracic vertebra. It passes beneath the aortic arch, crosses in front of the esophagus, the thoracic duct, and the descending aorta, and has the left pulmonary artery lying at first above, and then in front of it. The left bronchus has no eparterial branch, and therefore it has been supposed by some that there is no upper lobe to the left lung, but that the so-called upper lobe corresponds to the middle lobe of the right lung. The left main bronchus divides into two secondary bronchi or lobar bronchi, to deliver air to the two lobes of the left lung—the superior and the inferior lobe.
The secondary bronchi divide further into tertiary bronchi, (also known as segmental bronchi), each of which supplies a bronchopulmonary segment. A bronchopulmonary segment is a division of a lung separated from the rest of the lung by a septum of connective tissue. This property allows a bronchopulmonary segment to be surgically removed without affecting other segments. Initially, there are ten segments in each lung, but during development with the left lung having just two lobes, two pairs of segments fuse to give eight, four for each lobe. The tertiary bronchi divide further in another three branchings known as 4th order, 5th order and 6th order segmental bronchi which are also referred to as subsegmental bronchi. These branch into many smaller bronchioles which divide into terminal bronchioles, each of which then gives rise to several respiratory bronchioles, which go on to divide into two to eleven alveolar ducts. There are five or six alveolar sacs associated with each alveolar duct. The alveolus is the basic anatomical unit of gas exchange in the lung.
The main bronchi have relatively large lumens that are lined by respiratory epithelium. This cellular lining has cilia departing towards the mouth which removes dust and other small particles. There is a smooth muscle layer below the epithelium arranged as two ribbons of muscle that spiral in opposite directions. This smooth muscle layer contains seromucous glands, which secrete mucus, in its wall. Hyaline cartilage is present in the bronchi, surrounding the smooth muscle layer. In the main bronchi, the cartilage forms C-shaped rings like those in the trachea, while in the smaller bronchi, hyaline cartilage is present in irregularly arranged crescent-shaped plates and islands. These plates give structural support to the bronchi and keep the airway open.
The bronchial wall normally has a thickness of 10% to 20% of the total bronchial diameter.
Microanatomy
The cartilage and mucous membrane of the main bronchus (primary bronchi) are similar to those in the trachea. They are lined with respiratory epithelium, which is classified as ciliated pseudostratified columnar epithelium. The epithelium in the main bronchi contains goblet cells, which are glandular, modified simple columnar epithelial cells that produce mucins, the main component of mucus. Mucus plays an important role in keeping the airways clear in the mucociliary clearance process.
As branching continues through the bronchial tree, the amount of hyaline cartilage in the walls decreases until it is absent in the bronchioles. As the cartilage decreases, the amount of smooth muscle increases. The mucous membrane also undergoes a transition from ciliated pseudostratified columnar epithelium, to simple ciliated cuboidal epithelium, to simple squamous epithelium in the alveolar ducts and alveoli
Variation
In 0.1 to 5% of people there is a right superior lobe bronchus arising from the main stem bronchus prior to the carina. This is known as a tracheal bronchus, and seen as an anatomical variation. It can have multiple variations and, although usually asymptomatic, it can be the root cause of pulmonary disease such as a recurrent infection. In such cases resection is often curative.
The cardiac bronchus has a prevalence of ≈0.3% and presents as an accessory bronchus arising from the bronchus intermedius between the upper lobar bronchus and the origin of the middle and lower lobar bronchi of the right main bronchus.
An accessory cardiac bronchus is usually an asymptomatic condition but may be associated with persistent infection or hemoptysis. In about half of observed cases the cardiac bronchus presents as a short dead-ending bronchial stump, in the remainder the bronchus may exhibit branching and associated aerated lung parenchyma.
Function
The bronchi function to carry air that is breathed in through to the functional tissues of the lungs, called alveoli. Exchange of gases between the air in the lungs and the blood in the capillaries occurs across the walls of the alveolar ducts and alveoli. The alveolar ducts and alveoli consist primarily of simple squamous epithelium, which permits rapid diffusion of oxygen and carbon dioxide.
Clinical significance
Bronchial wall thickening, as can be seen on CT scan, generally (but not always) implies inflammation of the bronchi (bronchitis). Normally, the ratio of the bronchial wall thickness and the bronchial diameter is between 0.17 and 0.23.
Bronchitis
Bronchitis is defined as inflammation of the bronchi, which can either be acute or chronic. Acute bronchitis is usually caused by viral or bacterial infections. Many sufferers of chronic bronchitis also suffer from chronic obstructive pulmonary disease (COPD), and this is usually associated with smoking or long-term exposure to irritants.
Aspiration
The left main bronchus departs from the trachea at a greater angle than that of the right main bronchus. The right bronchus is also wider than the left and these differences predispose the right lung to aspirational problems. If food, liquids, or foreign bodies are aspirated, they will tend to lodge in the right main bronchus. Bacterial pneumonia and aspiration pneumonia may result.
If a tracheal tube used for intubation is inserted too far, it will usually lodge in the right bronchus, allowing ventilation only of the right lung.
Asthma
Asthma is marked by hyperresponsiveness of the bronchi with an inflammatory component, often in response to allergens.
In asthma, the constriction of the bronchi can result in difficulty in breathing giving shortness of breath; this can lead to a lack of oxygen reaching the body for cellular processes. In this case, an inhaler can be used to rectify the problem. The inhaler administers a bronchodilator, which serves to soothe the constricted bronchi and to re-expand the airways. This effect occurs quite quickly.
Bronchial atresia
Bronchial atresia is a rare congenital disorder that can have a varied appearance. A bronchial atresia is a defect in the development of the bronchi, affecting one or more bronchi – usually segmental bronchi and sometimes lobar. The defect takes the form of a blind-ended bronchus. The surrounding tissue secretes mucus normally but builds up and becomes distended. This can lead to regional emphysema.
The collected mucus may form a mucoid impaction or a bronchocele, or both. A pectus excavatum may accompany a bronchial atresia.
Additional images
Citations
Sources
Moore, Keith L. and Arthur F. Dalley. Clinically Oriented Anatomy, 4th ed. (1999). .
External links
Respiratory system anatomy
Thorax (human anatomy) | Biology and health sciences | Respiratory system | Biology |
308680 | https://en.wikipedia.org/wiki/Amphisbaenia | Amphisbaenia | Amphisbaenia (called amphisbaenians or worm lizards) is a group of typically legless lizards, comprising over 200 extant species. Amphisbaenians are characterized by their long bodies, the reduction or loss of the limbs, and rudimentary eyes. As many species have a pink body and scales arranged in rings, they have a superficial resemblance to earthworms. While the genus Bipes retains forelimbs, all other genera are limbless. Phylogenetic studies suggest that they are nested within Lacertoidea, closely related to the lizard family Lacertidae. Amphisbaenians are widely distributed, occurring in North America, Europe, Africa, South America, Western Asia and the Caribbean. Most species are less than long.
Description
Despite a superficial resemblance to some primitive snakes, amphisbaenians have many unique features that distinguish them from other reptiles. Internally, their right lung is reduced in size to fit their narrow bodies, whereas in snakes, it is always the left lung. Their skeletal structure and skin are also different from those of other squamates. Both genetic and recent fossil evidence indicate that amphisbaenians lost their legs independently from snakes.
The head is stout, not set off from the neck, and either rounded, sloped, or sloped with a ridge down the middle. Most of the skull is solid bone, with a distinctive single median tooth in the upper jaw. It has no outer ears, and the eyes are deeply recessed and covered with skin and scales. These rudimentary eyes have a cornea, lens, and complex ciliary body, which allows them to detect light, but they are reduced in size and do not have an anterior chamber. The body is elongated, and the tail truncates in a manner that vaguely resembles the head. At their tail is a single fracture plane for tail autotomy, between the fifth and eighth caudal rings and is often visible due to coloration. The purpose seems to be to distract predators with the tail acting as a decoy. Their name is derived from Amphisbaena, a mythical serpent with a head at each end—referencing both the manner in which their tail truncates, and their ability to move just as well in reverse as forwards. The four species of Bipes are unusual in having a pair of forelimbs. All other species lack any trace of forelimb skeletal elements, and Rhineura floridana also lack any pectoral girdle skeletal element. The other species have some remnants of the pectoral girdle embedded within the body musculature. A remnant of the pelvic girdle is present in all families, and Bipes and the genus Blanus have also retained a reduced femur.
Amphisbaenians have a distinctive skin made up of rings of scales (annuli) that form a tube in which the loosely attached trunk of the body moves. Burrowing is achieved with an accordion-like motion, with longitudinal muscles in the skin bunching up the annuli, anchoring it to the surrounding soil, and trunk muscles moving the body forward or backwards within the integumentary tube.
Amphisbaenians are carnivorous, able to tear chunks out of larger prey with their powerful, interlocking teeth. Like lizards, some species are able to shed their tails (autotomy). Most species lay eggs, although at least some are known to be viviparous.
The red worm lizard (Amphisbaena alba) is often found in association with leafcutter ants. This reptile is thought to forage in the ants' deep galleries, where the insects deposit their waste. The presence of these reptiles is easily explained by the fact that they prey on the larvae of large beetles that also inhabit the leafcutter ants' galleries.
Amphisbaenians have often been categorized by their skull shape. The specialized skull shape is hypothesized to be driven by environmental and ecological conditions, such as soil type, and is an instance of convergent evolution. Traditionally four types of skulls are recognized; “shovel-headed,” “round-headed,” “keel-headed,” and “spade-headed”, although it doesn't say anything about the relationship between the types. Of these four morphotypes, the round-headed species produce the lowest burrowing forces, the shovel-headed species the second lowest forces, the keel-headed species the second highest forces, and the spade-headed the highest forces.
Distribution
Amphisbaenians are found in North America, Europe, Africa, South America, the Middle East, and the Caribbean, a surprisingly large distribution despite being small subterranean animals that rarely ever leave their burrows.
Initially, this large distribution was thought to be due to vicariance, or the result of the breakup of Pangaea. This hypothesis was supported by morphological data that dated amphisbaenian diversification to over 200 million years ago (Mya), while Pangaea was still intact. However, a recent study using a combination of molecular and fossil evidence suggests that amphisbaenians originated in North America, where they underwent their first divergence around 107 Mya. They then underwent another major diversification into North American and European forms 40–56 Mya. Finally, the African and South American forms split around 40 Mya. This suggests that worm-lizards crossed the Atlantic Ocean (which had fully formed by 100 Mya) twice, once just after the K–Pg extinction, and then again, later in the Palaeogene. This also implies that limblessness evolved independently three times, a finding that contrasts the morphological theory that limbed amphisbaenians are the most basal. This widespread dispersal is suggested to have occurred by rafting – natural erosion or a storm event loosened a large raft of soil and vegetation that drifted across the ocean until landing on another shore. This oceanic rafting would be feasible due to the subterranean lifestyle and small nutritional requirements of amphisbaenids. After the Chicxulub impact, many predators of amphisbaenians became extinct, which allowed colonist amphisbaenians to thrive in new territories.
Evolution
The fully limbed Slavoia darevskii from the Late Cretaceous (Campanian) of Mongolia may represent an early relative of amphisbaenians. The oldest known modern amphisbaenians include members of Rhineuridae and the extinct family Oligodontosauridae from the Paleocene of North America. Modern amphisbaenians likely originated in North America, before dispersing to South America, Africa and Europe via rafting during the Paleogene.
Taxonomy
Taxonomic classification of amphisbaenians was traditionally based on morphological characters, such as the number of preanal pores, body annuli, tail annuli, and skull shape. Such characters are vulnerable to convergent evolution; in particular, the loss of the forelimbs and the evolution of specialized shovel-headed and keel-headed morphs appear to have occurred multiple times in the history of the group. Classifications based on mitochondrial DNA sequences and nuclear DNA sequences better reflect their true evolutionary history, and are now being used to distinguish genera of amphisbaenians.
The most ancient branch of the tree is the Rhineuridae. The remaining five families form a group to the exclusion of rhineurids. Bipedidae, Blanidae, and Cadeidae represent the most ancient divergences within this grouping, with Trogonophidae and Amphisbaenidae diverging more recently. South American amphisbaenids apparently are derived from African amphisbaenids that rafted across the Atlantic in the Eocene, about 40 million years ago. Cuban cadeids may be similarly derived from blanids that rafted across from northwestern Africa or southwestern Europe in a similar time frame.
Historically considered to be lizards, some studies have suggested that they should be considered separate from lizards, though many modern studies consider them to be true lizards, as they are closely related to other lizards of the clade Lacertoidea.
Families
Six extant families of amphisbaenians are currently recognised:
Amphisbaenidae Gray, 1865 – Amphisbaenids, tropical worm lizards of South America, some Caribbean islands, and Sub-Saharan Africa. (12 genera, 182 species)
Bipedidae Taylor, 1951 – Only in Mexico and commonly called ajolotes, but not to be confused with axolotls (1 genus, 3 species)
Blanidae Kearney & Stuart, 2004 – Anatolian, Iberian, and Moroccan worm lizards (1 genus, 7 species)
Cadeidae Vidal and Hedges, 2008 – Cuban keel-headed worm lizards (1 genus, 2 species). Traditionally amphisbaenids, but shown by DNA to be closest to Blanidae.
Rhineuridae Vanzolini, 1951 – North American worm lizards (1 genus, 1 species)
Trogonophidae Gray, 1865 – Palearctic worm lizards (4 genera, 6 species)
In addition, the following extinct families are also known from fossil remains:
†Chthonophidae Longrich et al., 2015
†Oligodontosauridae Estes, 1975
†Polyodontobaenidae Folie, Smith & Smith, 2013
Another fossil family, the †Crythiosauridae, was also previously placed in this group, but has since been removed due to a lack of evidence placing it within the amphisbaenians.
Phylogeny
The following cladogram shows the relationships between the six amphisbaenian families determined in the phylogenetic analysis of mitochondrial and nuclear genes by Vidal et al. (2008).
| Biology and health sciences | Lizards and other Squamata | Animals |
308865 | https://en.wikipedia.org/wiki/Polymer%20backbone | Polymer backbone | In polymer science, the polymer chain or simply backbone of a polymer is the main chain of a polymer. Polymers are often classified according to the elements in the main chains. The character of the backbone, i.e. its flexibility, determines the properties of the polymer (such as the glass transition temperature). For example, in polysiloxanes (silicone), the backbone chain is very flexible, which results in a very low glass transition temperature of . The polymers with rigid backbones are prone to crystallization (e.g. polythiophenes) in thin films and in solution. Crystallization in its turn affects the optical properties of the polymers, its optical band gap and electronic levels.
Organic polymers
Common synthetic polymers have main chains composed of carbon, i.e. C-C-C-C.... Examples include polyolefins such as polyethylene ((CH2CH2)n) and many substituted derivative ((CH2CH(R))n) such as polystyrene (R = C6H5), polypropylene (R = CH3), and acrylates (R = CO2R').
Other major classes of organic polymers are polyesters and polyamides. They have respectively -C(O)-O- and -C(O)-NH- groups in their backbones in addition to chains of carbon. Major commercial products are polyethyleneterephthalate ("PET"), ((C6H4CO2C2H4OC(O))n) and nylon-6 ((NH(CH2)5C(O))n).
Inorganic polymers
Siloxanes are a premier example of an inorganic polymer, even though they have extensive organic substituents. Their backbond is composed of alternating silicon and oxygen atoms, i.e. Si-O-Si-O... The silicon atoms bear two substituents, usually methyl as in the case of polydimethylsiloxane. Some uncommon but illustrative inorganic polymers include polythiazyl ((SN)x) with alternating S and N atoms, and polyphosphates ((PO3−)n).
Biopolymers
Major families of biopolymers are polysaccharides (carbohydrates), peptides, and polynucleotides. Many variants of each are known.
Proteins and peptides
Proteins are characterized by amide linkages (-N(H)-C(O)-) formed by the condensation of amino acids. The sequence of the amino acids in the polypeptide backbone is known as the primary structure of the protein. Like almost all polymers, protein fold and twist, forming into the secondary structure, which is rigidified by hydrogen bonding between the carbonyl oxygens and amide hydrogens in the backbone, i.e. C=O---HN. Further interactions between residues of the individual amino acids form the protein's tertiary structure. For this reason, the primary structure of the amino acids in the polypeptide backbone is the map of the final structure of a protein, and it therefore indicates its biological function. Spatial positions of backbone atoms can be reconstructed from the positions of alpha carbons using computational tools for the backbone reconstruction.
Carbohydrates
Carbohydrates arise by condensation of monosaccharides such as glucose. The polymers can be classified into oligosaccharides (up to 10 residues) and polysaccharides (up to about 50,000 residues). The backbone chain is characterized by an ether bond between individual monosaccharides. This bond is called the glycosidic linkage. These backbone chains can be unbranched (containing one linear chain) or branched (containing multiple chains). The glycosidic linkages are designated as alpha or beta depending on the relative stereochemistry of the anomeric (or most oxidized) carbon. In a Fischer Projection, if the glycosidic linkage is on the same side or face as carbon 6 of a common biological saccharide, the carbohydrate is designated as beta and if the linkage is on the opposite side it is designated as alpha. In a traditional "chair structure" projection, if the linkage is on the same plane (equatorial or axial) as carbon 6 it is designated as beta and on the opposite plane it is designated as alpha. This is exemplified in sucrose (table sugar) which contains a linkage that is alpha to glucose and beta to fructose. Generally, carbohydrates which our bodies break down are alpha-linked (example: glycogen) and those which have structural function are beta-linked (example: cellulose).
Nucleic acids
Deoxyribonucleic acid (DNA) and ribonucleic acid (RNA) are the main examples of polynucleotides. They arise by condensation of nucleotides. Their backbones form by the condensation of a hydroxy group on a ribose with the phosphate group on another ribose. This linkage is called a phosphodiester bond. The condensation is catalyzed by enzymes called polymerases. DNA and RNA can be millions of nucleotides long thus allowing for the genetic diversity of life. The bases project from the pentose-phosphate polymer backbone and are hydrogen bonded in pairs to their complementary partners (A with T and G with C). This creates a double helix with pentose phosphate backbones on either side, thus forming a secondary structure.
| Physical sciences | Concepts_2 | Chemistry |
308870 | https://en.wikipedia.org/wiki/Ring-opening%20polymerization | Ring-opening polymerization | In polymer chemistry, ring-opening polymerization (ROP) is a form of chain-growth polymerization in which the terminus of a polymer chain attacks cyclic monomers to form a longer polymer (see figure). The reactive center can be radical, anionic or cationic.
Ring-opening of cyclic monomers is often driven by the relief of bond-angle strain. Thus, as is the case for other types of polymerization, the enthalpy change in ring-opening is negative. Many rings undergo ROP.
Monomers
Many cyclic monomers are amenable to ROP. These include epoxides, cyclic trisiloxanes, some lactones and lactides, cyclic anhydrides, cyclic carbonates, and amino acid N-carboxyanhydrides. Many strained cycloalkenes, e.g norbornene, are suitable monomers via ring-opening metathesis polymerization. Even highly strained cycloalkane rings, such as cyclopropane and cyclobutane derivatives, can undergo ROP.
History
Ring-opening polymerization has been used since the beginning of the 1900s to produce polymers. Synthesis of polypeptides which has the oldest history of ROP, dates back to the work in 1906 by Leuchs. Subsequently, the ROP of anhydro sugars provided polysaccharides, including synthetic dextran, xanthan gum, welan gum, gellan gum, diutan gum, and pullulan. Mechanisms and thermodynamics of ring-opening polymerization were established in the 1950s. The first high-molecular weight polymers (Mn up to 105) with a repeating unit were prepared by ROP as early as in 1976.
An industrial application is the production of nylon-6 from caprolactam.
Mechanisms
Ring-opening polymerization can proceed via radical, anionic, or cationic polymerization as described below. Additionally, radical ROP is useful in producing polymers with functional groups incorporated in the backbone chain that cannot otherwise be synthesized via conventional chain-growth polymerization of vinyl monomers. For instance, radical ROP can produce polymers with ethers, esters, amides, and carbonates as functional groups along the main chain.
Anionic ring-opening polymerization (AROP)
Anionic ring-opening polymerizations (AROP) involve nucleophilic reagents as initiators. Monomers with a three-member ring structure - such as epoxides, aziridines, and episulfides - undergo anionic ROP.
A typical example of anionic ROP is that of ε-caprolactone, initiated by an alkoxide.
Cationic ring-opening polymerization
Cationic initiators and intermediates characterize cationic ring-opening polymerization (CROP). Examples of cyclic monomers that polymerize through this mechanism include lactones, lactams, amines, and ethers. CROP proceeds through an SN1 or SN2 propagation, chain-growth process. The mechanism is affected by the stability of the resulting cationic species. For example, if the atom bearing the positive charge is stabilized by electron-donating groups, polymerization will proceed by the SN1 mechanism. The cationic species is a heteroatom and the chain grows by the addition of cyclic monomers thereby opening the ring system.
The monomers can be activated by Bronsted acids, carbenium ions, onium ions, and metal cations.
CROP can be a living polymerization and can be terminated by nucleophilic reagents such as phenoxy anions, phosphines, or polyanions. When the amount of monomers becomes depleted, termination can occur intra or intermolecularly. The active end can "backbite" the chain, forming a macrocycle. Alkyl chain transfer is also possible, where the active end is quenched by transferring an alkyl chain to another polymer.
Ring-opening metathesis polymerization
Ring-opening metathesis polymerisation (ROMP) produces unsaturated polymers from cycloalkenes or bicycloalkenes. It requires organometallic catalysts.
The mechanism for ROMP follows similar pathways as olefin metathesis. The initiation process involves the coordination of the cycloalkene monomer to the metal alkylidene complex, followed by a [2+2] type cycloaddition to form the metallacyclobutane intermediate that cycloreverts to form a new alkylidene species.
Commercially relevant unsaturated polymers synthesized by ROMP include polynorbornene, polycyclooctene, and polycyclopentadiene.
Thermodynamics
The formal thermodynamic criterion of a given monomer polymerizability is related to a sign of the free enthalpy (Gibbs free energy) of polymerization:
where:
and indicate monomer and polymer states, respectively ( and/or = l (liquid), g (gaseous), c (amorphous solid), c' (crystalline solid), s (solution));
is the enthalpy of polymerization (SI unit: joule per kelvin);
is the entropy of polymerization (SI unit: joule);
is the absolute temperature (SI unit: kelvin).
The free enthalpy of polymerization () may be expressed as a sum of standard enthalpy of polymerization () and a term related to instantaneous monomer molecules and growing macromolecules concentrations:
where:
is the gas constant;
is the monomer;
is the monomer in an initial state;
is the active monomer.
Following Flory–Huggins solution theory that the reactivity of an active center, located at a macromolecule of a sufficiently long macromolecular chain, does not depend on its degree of polymerization (), and taking in to account that (where and indicate a standard polymerization enthalpy and entropy, respectively), we obtain:
At equilibrium (), when polymerization is complete the monomer concentration () assumes a value determined by standard polymerization parameters ( and ) and polymerization temperature:
Polymerization is possible only when . Eventually, at or above the so-called ceiling temperature (), at which , formation of the high polymer does not occur.
For example, tetrahydrofuran (THF) cannot be polymerized above = 84 °C, nor cyclo-octasulfur (S8) below = 159 °C. However, for many monomers, and , for polymerization in the bulk, are well above or below the operable polymerization temperatures, respectively.
The polymerization of a majority of monomers is accompanied by an entropy decrease, due mostly to the loss in the translational degrees of freedom. In this situation, polymerization is thermodynamically allowed only when the enthalpic contribution into prevails (thus, when and , the inequality is required). Therefore, the higher the ring strain, the lower the resulting monomer concentration at equilibrium.
Additional reading
| Physical sciences | Organic reactions | Chemistry |
309089 | https://en.wikipedia.org/wiki/Dependent%20personality%20disorder | Dependent personality disorder | Dependent personality disorder (DPD) is a personality disorder characterized by a pervasive psychological dependence on other people. This personality disorder is a long-term condition in which people depend on others to meet their emotional and physical needs. Dependent personality disorder is a cluster C personality disorder, which is characterized by excessive fear and anxiety. It begins prior to early adulthood, and it is present in a variety of contexts and is associated with inadequate functioning. Symptoms can include anything from extreme passivity, devastation or helplessness when relationships end, avoidance of responsibilities, and severe submission.
Signs and symptoms
People who have dependent personality disorder are overdependent on other people when it comes to making decisions. They cannot make a decision on their own as they need constant approval from other people. Consequently, individuals diagnosed with DPD tend to place needs and opinions of others above their own as they do not have the confidence to trust their decisions. This kind of behavior can explain why people with DPD tend to show passive and clingy behaviour. These individuals display a fear of separation and cannot stand being alone. When alone, they experience feelings of isolation and loneliness due to their overwhelming dependence on other people. Generally people with DPD are also pessimistic: they expect the worst out of situations or believe that the worst will happen. They tend to be more introverted and are more sensitive to criticism and fear rejection.
Risk factors
People with a history of neglect and an abusive upbringing are more susceptible to develop DPD, specifically those involved in long-term abusive relationships. Those with overprotective or authoritarian parents are also more at risk to develop DPD. Having a family history of anxiety disorder can play a role in the development of DPD as a 2004 twin study found a 0.81 heritability for personality disorders collectively.
Causes
While the exact cause of dependent personality disorder is unknown, a study in 2012 estimated that between 55% and 72% of the risk of the condition is inherited from one's parents. The difference between a "dependent personality" and a "dependent personality disorder" is somewhat subjective, which makes diagnosis sensitive to cultural influences such as gender role expectations.
Dependent traits in children tended to increase with parenting behaviours and attitudes characterized by overprotectiveness and authoritarianism. Thus the likelihood of developing dependent personality disorder increased, since these parenting traits can limit them from developing a sense of autonomy, rather teaching them that others are powerful and competent.
Traumatic or adverse experiences early in an individual's life, such as neglect and abuse or serious illness, can increase the likelihood of developing personality disorders, including dependent personality disorder, later on in life. This is especially prevalent for those individuals who also experience high interpersonal stress and poor social support.
There is a higher frequency of the disorder seen in women than men; hence, expectations relating to gender role may contribute to some extent.
Diagnosis
Clinicians and clinical researchers conceptualize dependent personality disorder in terms of four related components:
Cognitive: a perception of oneself as powerless and ineffectual, coupled with the belief that other people are comparatively powerful and potent.
Motivational: a desire to obtain and maintain relationships with protectors and caregivers.
Behavioral: a pattern of relationship-facilitating behavior designed to strengthen interpersonal ties and minimize the possibility of abandonment and rejection.
Emotional: fear of abandonment, fear of rejection, and anxiety regarding evaluation by figures of authority.
American Psychiatric Association and DSM
The Diagnostic and Statistical Manual of Mental Disorders (DSM) contains a dependent personality disorder diagnosis. It refers to a pervasive and excessive need to be taken care of which leads to submissive and clinging behavior and fears of separation. This begins prior to early adulthood and can be present in a variety of contexts.
In the DSM Fifth Edition (DSM-5), there is one criterion by which there are eight features of dependent personality disorder. The disorder is indicated by at least five of the following factors:
Has difficulty making everyday decisions without an excessive amount of advice and reassurance from others.
Needs others to assume responsibility for most major areas of their life.
Has difficulty expressing disagreement with others because of fear of loss of support or approval.
Has difficulty initiating projects or doing things on their own (because of a lack of self confidence in judgment or abilities rather than a lack of motivation or energy).
Goes to excessive lengths to obtain nurturance and support from others, to the point of volunteering to do things that are unpleasant.
Feels uncomfortable or helpless when alone because of exaggerated fears of being unable to care for themselves.
Urgently seeks another relationship as a source of care and support when a close relationship ends.
Is unrealistically preoccupied with fears of being left to take care of themselves.
The diagnosis of personality disorders in the fourth edition Diagnostic and Statistical Manual of Mental Disorders, including dependent personality disorder, was found to be problematic due to reasons such as excessive diagnostic comorbidity, inadequate coverage, arbitrary boundaries with normal psychological functioning, and heterogeneity among individuals within the same categorial diagnosis.
World Health Organization
The World Health Organization's ICD-10 lists dependent personality disorder as Dependent personality disorder:
It is characterized by at least 4 of the following:
Encouraging or allowing others to make most of one's important life decisions;
Subordination of one's own needs to those of others on whom one is dependent, and undue compliance with their wishes;
Unwillingness to make even reasonable demands on the people one depends on;
Feeling uncomfortable or helpless when alone, because of exaggerated fears of inability to care for oneself;
Preoccupation with fears of being abandoned by a person with whom one has a close relationship, and of being left to care for oneself;
Limited capacity to make everyday decisions without an excessive amount of advice and reassurance from others.
Associated features may include perceiving oneself as helpless, incompetent, and lacking stamina.
Includes:
Asthenic, inadequate, passive, and self-defeating personality (disorder)
It is a requirement of ICD-10 that a diagnosis of any specific personality disorder also satisfies a set of general personality disorder criteria.
SWAP-200
The SWAP-200 is a diagnostic tool that was proposed with the goal of overcoming limitations, such as limited external validity for the diagnostic criteria for dependent personality disorder, to the DSM. It serves as a possible alternative nosological system that emerged from the efforts to create an empirically based approach to personality disorders – while also preserving the complexity of clinical reality. Dependent personality disorder is considered a clinical prototype in the context of the SWAP-200. Rather than discrete symptoms, it provides composite description characteristic criteria – such as personality tendencies.
Based on the Q-Sort method and prototype matching, the SWAP-200 is a personality assessment procedure relying on an external observer's judgment. It provides:
A personality diagnosis expressed as the matching with ten prototypical descriptions of DSM-IV personality disorders.
A personality diagnosis based on the matching of the patient with 11 Q-factors of personality derived empirically.
A dimensional profile of healthy and adaptive functioning.
The traits that define dependent personality disorder according to SWAP-200 are:
They tend to become attached quickly and/or intensely, developing feelings and expectations that are not warranted by the history or context of the relationship.
Since they tend to be ingratiating and submissive, people with DPD tend to be in relationships in which they are emotionally or physically abused.
They tend to feel ashamed, inadequate, and depressed.
They also feel powerless and tend to be suggestible.
They are often anxious and tend to feel guilty.
These people have difficulty acknowledging and expressing anger and struggle to get their own needs and goals met.
Unable to soothe or comfort themselves when distressed, they require involvement of another person to help regulate their emotions.
Psychodynamic Diagnostic Manual
The Psychodynamic Diagnostic Manual (PDM) approaches dependent personality disorder in a descriptive, rather than prescriptive sense and has received empirical support. The Psychodynamic Diagnostic Manual includes two different types of dependent personality disorder:
Passive-aggressive
Counter-dependent
The PDM-2 adopts and applies a prototypic approach, using empirical measures like the SWAP-200. It was influenced by a developmental and empirically grounded perspective, as proposed by Sidney Blatt. This model is of particular interest when focusing on dependent personality disorder, claiming that psychopathology comes from distortions of two main coordinates of psychological development:
The anaclitic/introjective dimension.
The relatedness/self-definition dimension.
The anaclitic personality organization in individuals exhibits difficulties in interpersonal relatedness, exhibiting the following behaviours:
Preoccupation with relationships
Fear of abandonment and of rejection
Seeking closeness and intimacy
Difficulty managing interpersonal boundaries
Tend to have an anxious-preoccupied attachment style.
Introjective personality style is associated with problems in self-definition.
Differential diagnosis
There are similarities between individuals with dependent personality disorder and individuals with borderline personality disorder, in that they both have a fear of abandonment. Those with dependent personality disorder do not necessarily exhibit impulsive behaviour or unstable affect experienced by those with borderline personality disorder, differentiating the two disorders.
Treatment
People who have DPD are generally treated with psychotherapy. The main goal of this therapy is to make the individual more independent and help them form healthy relationships with the people around them. This is done by improving their self-esteem and confidence.
Medication can be used to treat patients who suffer from depression or anxiety because of their DPD, but this does not treat the core problems caused by the disorder.
Epidemiology
Based on a recent survey of 43,093 Americans, 0.49% of adults meet diagnostic criteria for DPD (National Epidemiologic Survey on Alcohol and Related Conditions; NESARC; Grant et al., 2004).
Traits related to DPD, like most personality disorders, emerge in childhood or early adulthood. Findings from the NESArC study found that 18 to 29 year olds have a greater chance of developing DPD. DPD is more common among women compared to men as 0.6% of women have DPD compared to 0.4% of men.
A 2004 twin study suggests a heritability of 0.81 for developing dependent personality disorder. Because of this, there is significant evidence that this disorder runs in families.
Children and adolescents with a history of anxiety disorders and physical illnesses are more susceptible to acquiring this disorder.
Millon's subtypes
Psychologist Theodore Millon identified five adult subtypes of dependent personality disorder. Any individual dependent may exhibit none or one or more of the following:
History
The conceptualization of dependency, within classical psychoanalytic theory, is directly related to Sigmund Freud's oral psychosexual stage of development. Frustration or over-gratification was said to result in an oral fixation and in an oral type of character, characterized by feeling dependent on others for nurturing and by behaviors representative of the oral stage. Later psychoanalytic theories shifted the focus from a drive-based approach of dependency to the recognition of the importance of early relationships and establishing separation from these early caregivers, in which the exchanges between the caregiver and the child become internalized, and the nature of these interactions becomes part of the concepts of the self and of others.
| Biology and health sciences | Mental disorders | Health |
309208 | https://en.wikipedia.org/wiki/Aerobic%20exercise | Aerobic exercise | Aerobic exercise, also known as cardio, is physical exercise of low to high intensity that depends primarily on the aerobic energy-generating process. "Aerobic" is defined as "relating to, involving, or requiring oxygen", and refers to the use of oxygen to meet energy demands during exercise via aerobic metabolism adequately. Aerobic exercise is performed by repeating sequences of light-to-moderate intensity activities for extended periods of time. According to the World Health Organization, over 31% of adults and 80% of adolescents fail to maintain the recommended levels of physical activity. Examples of cardiovascular or aerobic exercise are medium- to long-distance running or jogging, swimming, cycling, stair climbing and walking.
For reducing the risk of health issues, 2.5 hours of moderate-intensity aerobic exercise per week is recommended. At the same time, even doing an hour and a quarter (11 minutes/day) of exercise can reduce the risk of early death, cardiovascular disease, stroke, and cancer.
Aerobic exercise may be better referred to as "solely aerobic", as it is designed to be low-intensity enough that all carbohydrates are aerobically turned into energy via mitochondrial ATP production. Mitochondria are organelles that rely on oxygen for the metabolism of carbs, proteins, and fats. Aerobic exercise causes a remodeling of mitochondrial cells within the tissues of the liver and heart.
History
Archibald Hill, a British physiologist, introduced the concepts of maximal oxygen uptake and oxygen debt in 1922. German physician Otto Meyerhof and Hill shared the 1922 Nobel Prize in Physiology or Medicine for their independent work related to muscle energy metabolism. Building on this work, scientists began measuring oxygen consumption during exercise. Henry Taylor at the University of Minnesota and Swedish scientists Per-Olof Åstrand and Bengt Saltin made notable contributions in the 1950s and 60s. Contributions were also made by the Harvard Fatigue Laboratory, Copenhagen Muscle Research Centre as well as various German universities.
After World War II, health-oriented recreational activities such as jogging became popular. The Royal Canadian Air Force Exercise Plans, developed by Dr. Bill Orban and published in 1961, helped to launch modern fitness culture.
Physical therapists Col. Pauline Potts and Dr. Kenneth H. Cooper, both of the United States Air Force, advocated the concept of aerobic exercise. In the 1960s, Cooper started research into preventive medicine. He conducted the first extensive research on aerobic exercise on over 5,000 U.S. Air Force personnel after becoming intrigued by the belief that exercise can preserve one's health. In 1966 he coined the term "aerobics". Two years later, in 1968, he published a book of the same name. In 1970, he created the Cooper Institute for non-profit research and education devoted to preventive medicine. He published a mass-market version of his book The New Aerobics in 1979. Cooper encouraged millions into becoming active and is now known as the "father of aerobics". Cooper's book inspired Jacki Sorensen to create aerobic dancing exercise routines, which grew in popularity in the 1970s in the U.S., and at the same time, Judi Missett developed and expanded Jazzercise.
In the 1970s, there was a running boom. It was inspired by the Olympics, the New-York marathon and the advent of cushioned shoes.
Aerobics at home became popular worldwide after the release of Jane Fonda's Workout exercise video in 1982. Step aerobics was popular in the 1990s, driven by a step product and program from Reebok shoes.
Definition
Aerobic exercise comprises innumerable forms. In general, it is performed at a moderate level of intensity over a relatively long period of time. For example, running a long distance at a moderate pace is an aerobic exercise, but sprinting is not. Playing singles tennis, with near-continuous motion, is generally considered aerobic activity, while activities with brief bursts of energetic movement within longer periods of casual movement may not be aerobic. Some sports are thus inherently "aerobic", while other aerobic exercises, such as fartlek training or aerobic dance classes, are designed specifically to improve aerobic capacity and fitness. It is most common for aerobic exercises to involve the leg muscles, primarily or exclusively. There are some exceptions. For example, rowing to distances of 2,000 meters or more is an aerobic sport that exercises several major muscle groups, including those of the legs, abdominals, chest, and arms.
Examples
Moderate activities
Swimming
Dancing
Hiking on flat ground
Bicycling at less than
Moderate walking (about )
Downhill skiing
Tennis (doubles)
Softball
Gardening
Light yard work
Jogging
Vigorous activities
Brisk walking (about )
Bicycling at more than
Hiking uphill
Cross-country skiing
Stair climbing
Soccer
Jogging
Jumping rope
Tennis (singles)
Basketball
Heavy yard work
Elliptical training
Rowing
Versus anaerobic exercise
Aerobic exercise and fitness can be contrasted with anaerobic exercise, of which strength training and short-distance running are the most salient examples. The two types of exercise differ by the duration and intensity of muscular contractions involved, as well as by how energy is generated within the muscle. Common kettlebell exercises combine aerobic and anaerobic aspects. Allowing 24 hours of recovery between aerobic and strength exercise leads to greater fitness.
New research on the endocrine functions of contracting muscles has shown that both aerobic and anaerobic exercise promote the secretion of myokines, with attendant benefits including growth of new tissue, tissue repair, and various anti-inflammatory functions, which in turn reduce the risk of developing various inflammatory diseases. Myokine secretion in turn is dependent on the amount of muscle contracted, and the duration and intensity of contraction. As such, both types of exercise produce endocrine benefits.
In almost all conditions, anaerobic exercise is accompanied by aerobic (in the presence of oxygen) exercises because the less efficient anaerobic metabolism must supplement the aerobic system due to energy demands that exceed the aerobic system's capacity. During anaerobic exercise, the body must generate energy through other processes than aerobic metabolism, including glycolysis paired with lactic acid fermentation, and the phosphocreatine system to generate energy in the form of ATP.
Fuel usage
Depending on the intensity of exercise, the body preferentially utilizes certain fuel forms to meet energy demands. The two main fuel sources for aerobic exercise in the body include fat (in the form of adipose tissue) and glycogen. At lower intensity aerobic exercise, the body preferentially uses fat as its main fuel source for cellular respiration, however as intensity increases the body preferentially uses glycogen stored in the muscles and liver or other carbohydrates, as it is a quicker source of energy. Aerobic exercise at low or moderate intensity is not a very efficient way to lose fat in comparison to high intensity aerobic exercise. Lipolysis (hydrolysis of triglyceride into fatty acids), not fat burning (conversion of fatty acid to carbon dioxide), explains the intensity-dependent fat mass reduction. It has been shown that fatty acid is consumed for wound healing, where moderate intensity exercise does not produce significant damage like high intensity exercise. The size of adipose tissue is determined by the magnitude of nutrient competition from muscle and lungs for cell regeneration and energy replenishment after exercise.
Health effects
Among the possible health benefits of regular aerobic exercise are:
May improve mood
Strengthens and enlarges the heart muscle, to improve its pumping efficiency and reduce the resting heart rate, known as aerobic conditioning
May improve circulation efficiency and reduce blood pressure
May help maintain independence in later life
Increases the total number of red blood cells in the body, facilitating transport of oxygen
Improves mental health, including reducing stress and lowering the incidence of depression, as well as increased cognitive capacity.
Slightly reduced depression may also be observed, especially if aerobic exercises are used as additional treatment for patients with a hematological malignancy
Reduces the risk for diabetes (One meta-analysis has shown, from multiple conducted studies, that aerobic exercise does help lower Hb A1Clevels for type 2 diabetics.)
Moderates the risk of death due to cardiovascular problems
Promotes weight loss
Reduces the risk of osteoporosis
May improve episodic memory
Risks and disadvantages
Some drawbacks of aerobic exercise include:
Overuse injuries of the musculoskeletal system because of repetitive exercise, with young athletes (under the age of 19) particularly at risk
Overtraining syndrome may lead to persistent dysfunction of a number of body systems
High volumes of training with insufficient calorie intake puts athletes—particularly female ones—at risk for RED-S
Aerobic exercise may not be as efficient as other exercise methods. For example, High-intensity interval training (HIIT) has been shown to provide similar benefits in a fraction of the time spent exercising per week.
Both the health benefits and the performance benefits, or "training effect", require that the duration and the frequency of exercise both exceed a certain minimum. Most authorities suggest at least twenty minutes performed at least three times per week.
Commercialization
Aerobic exercise has long been a popular approach to achieving weight loss and physical fitness, often taking a commercial form.
In the 1970s, Judi Sheppard Missett helped create the market for commercial aerobics with her Jazzercise program, at the same time as Jacki Sorensen was expanding her system of aerobic dancing.
In the 1980s, Richard Simmons hosted an aerobic exercise show on television, and followed Jane Fonda's lead by releasing a series of exercise videos.
In the 1990s, Billy Blanks's Tae Bo helped popularize cardio-boxing workouts that incorporated martial arts movements. Reebok shoes popularized step aerobics with their Reebok Step device and training program.
| Biology and health sciences | Physical fitness | Health |
309252 | https://en.wikipedia.org/wiki/Phosphorescence | Phosphorescence | Phosphorescence is a type of photoluminescence related to fluorescence. When exposed to light (radiation) of a shorter wavelength, a phosphorescent substance will glow, absorbing the light and reemitting it at a longer wavelength. Unlike fluorescence, a phosphorescent material does not immediately reemit the radiation it absorbs. Instead, a phosphorescent material absorbs some of the radiation energy and reemits it for a much longer time after the radiation source is removed.
In a general sense, there is no distinct boundary between the emission times of fluorescence and phosphorescence (i.e.: if a substance glows under a black light it is generally considered fluorescent, and if it glows in the dark it is often simply called phosphorescent). In a modern, scientific sense, the phenomena can usually be classified by the three different mechanisms that produce the light, and the typical timescales during which those mechanisms emit light. Whereas fluorescent materials stop emitting light within nanoseconds (billionths of a second) after the excitation radiation is removed, phosphorescent materials may continue to emit an afterglow ranging from a few microseconds to many hours after the excitation is removed.
There are two separate mechanisms that may produce phosphorescence, called triplet phosphorescence (or simply phosphorescence) and persistent phosphorescence (or persistent luminescence):
Triplet phosphorescence occurs when an atom absorbs a high-energy photon, and the energy becomes locked in the spin multiplicity of the electrons, generally changing from a fluorescent singlet state to a slower emitting triplet state. The slower timescales of the reemission are associated with "forbidden" energy state transitions in quantum mechanics. As these transitions occur relatively slowly in certain materials, absorbed radiation is reemitted at a lower intensity, ranging from a few microseconds to as much as one second after the excitation is removed.
Persistent phosphorescence occurs when a high-energy photon is absorbed by an atom and its electron becomes trapped in a defect in the lattice of the crystalline or amorphous material. A defect such as a missing atom (vacancy defect) can trap an electron like a pitfall, storing that electron's energy until released by a random spike of thermal (vibrational) energy. Such a substance will then emit light of gradually decreasing intensity, ranging from a few seconds to up to several hours after the original excitation.
Everyday examples of phosphorescent materials are the glow-in-the-dark toys, stickers, paint, and clock dials that glow after being charged with a bright light such as in any normal reading or room light. Typically, the glow slowly fades out, sometimes within a few minutes or up to a few hours in a dark room.
The study of phosphorescent materials led to the discovery of radioactive decay.
Etymology
The term phosphorescence comes from the Ancient Greek word φῶς (phos), meaning "light", and the Greek suffix -φόρος (-phoros), meaning "to bear", combined with the Latin suffix -escentem, meaning "becoming of", "having a tendency towards", or "with the essence of". Thus, phosphorescence literally means "having a tendency to bear light". It was first recorded in 1766.
The term phosphor had been used since the Middle Ages to describe minerals that glowed in the dark. One of the most famous, but not the first, was Bolognian phosphor. Around 1604, Vincenzo Casciarolo discovered a "lapis solaris" near Bologna, Italy. Once heated in an oxygen-rich furnace, it thereafter absorbed sunlight and glowed in the dark. In 1677, Hennig Brand isolated a new element that glowed due to a chemiluminescent reaction when exposed to air, and named it "phosphorus".
In contrast, the term luminescence (from the Latin lumen for "light"), was coined by Eilhardt Wiedemann in 1888 as a term to refer to "light without heat", while "fluorescence" by Sir George Stokes in 1852, when he noticed that, when exposing a solution of quinine sulfate to light refracted through a prism, the solution glowed when exposed to the mysterious invisible-light (now known to be UV light) beyond the violet end of the spectrum. Stokes formed the term from a combination of fluorspar and opalescence (preferring to use a mineral instead of a solution), albeit it was later discovered that fluorspar glows due to phosphorescence.
There was much confusion between the meanings of these terms throughout the late nineteenth to mid-twentieth centuries. Whereas the term "fluorescence" tended to refer to luminescence that ceased immediately (by human-eye standards) when removed from excitation, "phosphorescence" referred to virtually any substance that glowed for appreciable periods in darkness, sometimes to include even chemiluminescence (which occasionally produced substantial amounts of heat). Only after the 1950s and 1960s did advances in quantum electronics, spectroscopy, and lasers provide a measure to distinguish between the various processes that emit the light, although in common speech the distinctions are still often rather vague.
Introduction
In simple terms, phosphorescence is a process in which energy absorbed by a substance is released relatively slowly in the form of light. This is in some cases the mechanism used for glow-in-the-dark materials which are "charged" by exposure to light. Unlike the relatively swift reactions in fluorescence, such as those seen in laser mediums like the common ruby, phosphorescent materials "store" absorbed energy for a longer time, as the processes required to reemit energy occur less often. However, timescale is still only a general distinction, as there are slow-emitting fluorescent materials, for example uranyl salts, and, likewise, some phosphorescent materials like zinc sulfide (in violet) are very fast. Scientifically, the phenomena are classified by the different mechanisms that produce the light, as materials that phosphoresce may be suitable for some purposes such as lighting, but may be completely unsuitable for others that require fluorescence, like lasers. Further blurring the lines, a substance may emit light by one, two, or all three mechanisms depending on the material and excitation conditions.
When the stored energy becomes locked in by the spin of the atomic electrons, a triplet state can occur, slowing the emission of light, sometimes by several orders of magnitude. Because the atoms usually begin in a singlet state of spin, favoring fluorescence, these types of phosphors typically produce both types of emission during illumination, and then a dimmer afterglow of strictly phosphorescent light typically lasting less than a second after the illumination is switched off.
Conversely, when the stored energy is due to persistent phosphorescence, an entirely different process occurs without a fluorescence precursor. When electrons become trapped within a defect in the atomic or molecular lattice, light is prevented from reemitting until the electron can escape. To escape, the electron needs a boost of thermal energy to help spring it out of the trap and back into orbit around the atom. Only then can the atom emit a photon. Thus, persistent phosphorescence is highly dependent on the temperature of the material.
Triplet phosphorescence
Most photoluminescent events, in which a chemical substrate absorbs and then re-emits a photon of light, are fast, in the order of 10 nanoseconds. Light is absorbed and emitted at these fast time scales in cases where the energy of the photons involved matches the available energy states and allowed transitions of the substrate. In the special case of phosphorescence, the electron which absorbed the photon (energy) undergoes an unusual intersystem crossing into an energy state of different (usually higher) spin multiplicity (see term symbol), usually a triplet state. As a result, the excited electron can become trapped in the triplet state with only "forbidden" transitions available to return to the lower energy singlet state. These transitions, although "forbidden", will still occur in quantum mechanics but are kinetically unfavored and thus progress at significantly slower time scales. Most phosphorescent compounds are still relatively fast emitters, with triplet decay-times in the order of milliseconds.
Common examples include the phosphor coatings used in fluorescent lamps, where phosphorescence on the order of milliseconds or longer is useful for filling in the "off-time" between AC current cycles, helping to reduce "flicker". Phosphors with faster decay times are used in applications like the pixels excited by free electrons (cathodoluminescence) in cathode-ray tube television-sets, which are slow enough to allow the formation of a picture as the electron beam scans the screen, but fast enough to prevent the frames from blurring together. Even substances commonly associated with fluorescence may in fact be prone to phosphorescence, such as the liquid dyes found in highlighter pens, which is a common problem in liquid dye lasers. The onset of phosphorescence in this case can sometimes be reduced or delayed significantly by the use of triplet-quenching agents.
Equation
where S is a singlet and T a triplet whose subscripts denote states (0 is the ground state, and 1 the excited state). Transitions can also occur to higher energy levels, but the first excited state is denoted for simplicity.
Persistent phosphorescence
Solid materials typically come in two main types: crystalline and amorphous. In either case, a lattice or network of atoms and molecules form. In crystals, the lattice is a very neat, uniform assembly. However, nearly all crystals have defects in the stacking sequence of these molecules and atoms. A vacancy defect, where an atom is simply missing from its place, leaving an empty "hole", is one type of defect. Sometimes atoms can move from place to place within the lattice, creating Schottky defects or Frenkel defects. Other defects can occur from impurities in the lattice. For example, when a normal atom is substituted by a different atom of much larger or smaller size, a substitutional defect occurs, while an interstitial defect occurs when a much smaller atom gets trapped in the "interstices", or the spaces between atoms. In contrast, amorphous materials have no "long-range order" (beyond the space of a few atoms in any direction), thus by definition are filled with defects.
When a defect occurs, depending on the type and material, it can create a hole, or a "trap". For example, a missing oxygen atom from a zinc oxide compound creates a hole in the lattice, surrounded by unbound zinc-atoms. This creates a net force or attraction that can be measured in electron-volts. When a high-energy photon strikes one of the zinc atoms, its electron absorbs the photon and is thrown out into a higher orbit. The electron may then enter the trap and be held in place (out of its normal orbit) by the attraction. To trigger the release of the energy, a random spike in thermal energy of sufficient magnitude is needed to boost the electron out of the trap and back into its normal orbit. Once in orbit, the electron's energy can drop back to normal (ground state) resulting in the release of a photon.
The release of energy in this way is a completely random process, governed mostly by the average temperature of the material versus the "depth" of the trap, or how many electron-volts it exerts. A trap that has a depth of 2.0 electron-volts would require a great amount of thermal energy (very high temperature) to overcome the attraction, while at a depth of 0.1 electron-volts very little heat (very cold temperature) is needed for the trap to even hold an electron. Generally, higher temperatures cause a faster release of energy, resulting in a brighter yet short-lived emission, while lower temperatures produce a dimmer but longer-lasting glow. Temperatures that are too hot or cold, depending on the substance, may not allow the accumulation or release of energy at all. The ideal depth of trap for persistent phosphorescence at room temperature is typically between 0.6 and 0.7 electron-volts. If the phosphorescent quantum yield is high, that is, if the substance has a large number of traps of the correct depth, this substance will release a significant amount of light over a long period of time, creating a so-called "glow in the dark" material.
Chemiluminescence
Some examples of glow-in-the-dark materials do not glow by phosphorescence. For example, glow sticks glow due to a chemiluminescent process which is commonly mistaken for phosphorescence. In chemiluminescence, an excited state is created via a chemical reaction. The light emission tracks the kinetic progress of the underlying chemical reaction. The excited state will then transfer to a dye molecule, also known as a sensitizer or fluorophor, and subsequently fluoresce back to the ground state.
Materials
Common pigments used in phosphorescent materials include zinc sulfide and strontium aluminate. Use of zinc sulfide for safety related products dates back to the 1930s.
The development of strontium aluminate pigments in 1993 was spurred on by the need to find a substitute for glow-in-the-dark materials with high luminance and long phosphorescence, especially those that used promethium. This led to the discovery by Yasumitsu Aoki (Nemoto & Co.) of materials with luminance approximately 10 times greater than zinc sulfide and phosphorescence approximately 10 times longer. This has relegated most zinc sulfide based products to the novelty category. Strontium aluminate based pigments are now used in exit signs, pathway marking, and other safety related signage.
Since both phosphorescence (transition from T1 to S0) and the generation of T1 from an excited singlet state (e.g., S1) via intersystem crossing (ISC) are spin-forbidden processes, most organic materials exhibit insignificant phosphorescence as they mostly fail to populate the excited triplet state, and, even if T1 is formed, phosphorescence is most frequently outcompeted by non-radiative pathways. One strategy to enhance the ISC and phosphorescence is the incorporation of heavy atoms, which increase spin-orbit coupling (SOC). Additionally, the SOC (and therefore the ISC) can be promoted by coupling n-π* and π-π* transitions with different angular momenta, also known as Mostafa El-Sayed's rule. Such transitions are typically exhibited by carbonyl or triazine derivatives, and most organic room-temperature phosphorescent (ORTP) materials incorporate such moieties. In turn, to inhibit competitive non-radiative deactivation pathways, including vibrational relaxation and oxygen quenching and triplet-triplet annihilations, organic phosphors have to be embedded in rigid matrices such as polymers, and molecular solids (crystals, covalent organic frameworks, and others).
Uses
In 1974 Becky Schroeder was given a US patent for her invention of the "Glow Sheet" which used phosphorescent lines under writing paper to help people write in low-light conditions.
Glow in the dark material is added to the plastic blend used in injection molds to make some disc golf discs, which allow the game to be played at night.
Often clock faces of watches are painted with phosphorescent colours. Therefore, they can be used in absolute dark environments for several hours after having been exposed to bright light.
A common use of phosphorescence is decoration. Stars made of glow-in-the-dark plastic are placed on walls, ceilings, or hanging from strings make a room look like the night sky. Other objects like figurines, cups, posters, lamp fixtures, toys and bracelet beads may also glow. Using blacklights makes these things glow brightly, common at raves, bedrooms, theme parks, and festivals.
Shadow wall
A shadow wall is created when a light flashes upon a person or object in front of a phosphorescent screen which temporarily captures the shadow. The screen or wall is painted with a glow-in-the-dark product that contains phosphorescent compounds. Publicly, these shadow walls can be found at certain science museums.
| Physical sciences | Electromagnetic radiation | Physics |
24169775 | https://en.wikipedia.org/wiki/Degradation%20%28geology%29 | Degradation (geology) | In geology, degradation refers to the lowering of a fluvial surface, such as a stream bed or floodplain, through erosional processes. Degradation is the opposite of aggradation. Degradation is characteristic of channel networks in which either bedrock erosion is taking place, or in systems that are sediment-starved and are therefore entraining more material than they are depositing. When a stream degrades, it leaves behind a fluvial terrace. This can be further classified as a strath terrace—a bedrock terrace that may have a thin mantle of alluvium—if the river is incising through bedrock. These terraces may be dated with methods such as cosmogenic radionuclide dating, OSL dating, and paleomagnetic dating (using reversals in the Earth's magnetic field to constrain the timing of events) to find when a river was at a particular level and how quickly it is downcutting.
| Physical sciences | Sedimentology | Earth science |
21193982 | https://en.wikipedia.org/wiki/Charcoal | Charcoal | Charcoal is a lightweight black carbon residue produced by strongly heating wood (or other animal and plant materials) in minimal oxygen to remove all water and volatile constituents. In the traditional version of this pyrolysis process, called charcoal burning, often by forming a charcoal kiln, the heat is supplied by burning part of the starting material itself, with a limited supply of oxygen. The material can also be heated in a closed retort. Modern charcoal briquettes used for outdoor cooking may contain many other additives, e.g. coal.
The history of wood charcoal production spans ancient times, rooted in the abundance of wood in various regions. The process typically involves stacking wood billets to form a conical pile, allowing air to enter through openings at the bottom, and igniting the pile gradually. Charcoal burners, skilled professionals tasked with managing the delicate operation, often lived in isolation to tend their wood piles. Throughout history, the extensive production of charcoal has been a significant contributor to deforestation, particularly in regions like Central Europe. However, various management practices, such as coppicing, aimed to maintain a steady supply of wood for charcoal production. The scarcity of easily accessible wood resources eventually led to the transition to fossil fuel equivalents like coal.
Modern methods of charcoal production involve carbonizing wood in retorts, yielding higher efficiencies compared to traditional kilning methods. The properties of charcoal depend on factors such as the material charred and the temperature of carbonization.
Charcoal finds diverse applications, including metallurgical fuel in iron and steel production, industrial fuel, cooking and heating fuel, reducing agent in chemical processes, and as a raw material in pyrotechnics. It is also utilized in cosmetics, horticulture, animal husbandry, medicine, and environmental sustainability efforts, such as carbon sequestration.
However, the production and utilization of charcoal can have adverse environmental impacts, including deforestation and emissions. Illegal and unregulated charcoal production, particularly in regions like South America and Africa, poses significant challenges to environmental conservation efforts.
History
Charcoal pile
The production of wood charcoal in locations where there is an abundance of wood dates back to ancient times. It generally began with piling billets of wood on their ends to form a conical pile. Openings were left at the bottom to admit air, with a central shaft serving as a flue. The whole pile was covered with turf or moistened clay. The firing began at the bottom of the flue, and the fire gradually spread outward and upward.
The traditional method in Britain used a charcoal pile or clamp. This was essentially a pile of wooden logs (e.g. seasoned oak) leaning in a circle against a chimney. The chimney consisted of 4 wooden stakes held up by some rope. In the clamp too the logs were completely covered with soil and straw allowing no air to enter. It must be lit by introducing some burning fuel into the chimney. The logs burned slowly and transformed into charcoal over a period of 5 days. If the soil covering became torn or cracked by the fire, additional soil was placed on the cracks. Once the burn was complete, the chimney was plugged to prevent air from entering.
Charcoal burners
The true art of this production method was in managing the sufficient generation of heat, by combusting part of the wood material, and the transfer of the heat to the wood in the process of being carbonized. The operation was so delicate that it was generally left to colliers (professional charcoal burners). They often lived alone in small huts to tend their wood piles. For example, in the Harz Mountains of Germany, charcoal burners lived in conical huts called Köten which still exist today.
Low efficiency and harmful emissions
The success of the operation depends upon the rate of the combustion. Under average conditions wood yields about 60% charcoal by volume, or 25% by weight; small-scale production methods often yield only about 50% by volume, while large-scale methods enabled higher yields of about 90% by the 17th century. A strong disadvantage of this production method is the huge amount of emissions that are harmful to human health and the environment (emissions of unburnt methane). As a result of the partial combustion of wood material, the efficiency of the traditional method is low.
Peak of production and decline
Deforestation and scarcity
The massive production of charcoal (at its height employing hundreds of thousands, mainly in Alpine and neighbouring forests) was a major cause of deforestation, especially in Central Europe. Complaints (as early as the Stuart period) about shortages may stem from over-exploitation or the impossibility of increasing production to match growing demand. In England, many woods were managed as coppices, which were cut and regrown cyclically, so that a steady supply of charcoal was available. But the increasing scarcity of easily harvested wood was a major factor behind the switch to fossil fuel equivalents, mainly coal and brown coal for industrial use.
By-product of wood tar production
In Finland and Scandinavia, the charcoal was considered the by-product of wood tar production. The best tar came from pine, thus pinewoods were cut down for tar pyrolysis. The residual charcoal was widely used as substitute for metallurgical coke in blast furnaces for smelting. Tar production led to rapid local deforestation. The end of tar production at the end of the 19th century resulted in rapid re-forestation of affected areas.
Charcoal briquette
The American form of the charcoal briquette was first invented and patented by Ellsworth B. A. Zwoyer of Pennsylvania in 1897 and was produced by the Zwoyer Fuel Company. The process was further popularized by Henry Ford, who used wood and sawdust byproducts from automobile fabrication as a feedstock. Ford Charcoal went on to become the Kingsford Company.
Production methods
The modern process of carbonizing wood, either in small pieces or as sawdust in cast iron retorts, is extensively practiced where wood is scarce, and also for the recovery of valuable byproducts (wood spirit, pyroligneous acid, wood tar), which the process permits. The question of the temperature of the carbonization is important; according to J. Percy, wood becomes brown at , a deep brown-black after some time at , and an easily powdered mass at . Charcoal made at is brown, soft and friable, and readily inflames at ; made at higher temperatures it is hard and brittle, and does not fire until heated to about .
Modern methods employ retorting technology, in which process heat is recovered from, and solely provided by, the combustion of gas released during carbonization. Yields of retorting are considerably higher than those of kilning, and may reach 35%-40%.
The properties of the charcoal produced depend on the material charred. The charring temperature is also important. Charcoal contains varying amounts of hydrogen and oxygen as well as ash and other impurities that, together with the structure, determine the properties. The approximate composition of charcoal for gunpowders is sometimes empirically described as C7H4O. To obtain a coal with high purity, source material should be free of non-volatile compounds.
Wood charcoal is obtained as the residue by destructive distillation of wood such that the products are:
Liquid products – pyroligneous acid and wood tar
Gaseous products – wood gas
Residual product – wood charcoal
Types
Common charcoal is made from peat, coal, wood, coconut shell, or petroleum.
Sugar charcoal is obtained from the carbonization of sugar and is particularly pure. It is purified by boiling with acids to remove any mineral matter and is then burned for a long time in a current of chlorine to remove the last traces of hydrogen. It was used by Henri Moissan in his early attempt to create synthetic diamonds.
Activated charcoal is similar to common charcoal but is manufactured especially for medical use. To produce activated charcoal, common charcoal is heated to about in the presence of an inert gas (usually argon or nitrogen), causing the charcoal to develop many internal spaces, or "pores", which help the activated charcoal to trap chemicals. Impurities on the surface of the charcoal are also removed during this process, greatly increasing its adsorption capacity.
Lump charcoal is a traditional charcoal made directly from hardwood material. It usually produces far less ash than briquettes.
Japanese charcoal has had pyroligneous acid removed during the charcoal making; it therefore produces almost no smell or smoke when burned. The traditional charcoal of Japan is classified into three types:
White charcoal (Binchōtan) is hard and produces a metallic sound when struck.
Ogatan is a more recent type made from hardened sawdust.
Pillow shaped briquettes are made by compressing charcoal, typically made from sawdust and other wood by-products, with a binder and other additives. The binder is usually starch. Briquettes may also include brown coal (heat source), mineral carbon (heat source), borax, sodium nitrate (ignition aid), limestone (ash-whitening agent), raw sawdust (ignition aid), and other additives.
Sawdust briquette charcoal is made by compressing sawdust without binders or additives. It is the preferred charcoal in Taiwan, Korea, Greece, and the Middle East. It has a round hole through the center, with a hexagonal cross-section. It is used primarily for barbecue as it produces no odor, no smoke, little ash, high heat, and has a long burning time (exceeding 4 hours).
Extruded charcoal is made by extruding either raw ground wood or carbonized wood into logs without the use of a binder. The heat and pressure of the extruding process hold the charcoal together. If the extrusion is made from raw wood material, the extruded logs are subsequently carbonized.
Uses
Charcoal has been used since earliest times for a large range of purposes including art and medicine, but by far its most important use has been as a metallurgical fuel. Charcoal is the traditional fuel of a blacksmith's forge and other applications where an intense heat is required. Charcoal was also used historically as a source of black pigment by grinding it up. In this form charcoal was important to early chemists and was a constituent of formulas for mixtures such as black powder. Due to its high surface area, charcoal can be used as a filter, catalyst, or adsorbent.
Metallurgical fuel
Charcoal burns at temperatures exceeding . By comparison, the melting point of iron is approximately . Due to its porosity, it is sensitive to the flow of air and the heat generated can be moderated by controlling the air flow to the fire. For this reason charcoal is still widely used by blacksmiths. Charcoal has been used for the production of iron and steel (where it also provided the necessary carbon) since at least 2000 BCE, with artifacts having been found in Proto-Hittite layers at Kaman-Kalehöyük. Charcoal briquettes can burn up to approximately with a forced air blower forge.
In the 16th century, England had to pass laws to prevent the country from becoming completely denuded of trees due to production of iron. In the 19th century charcoal was largely replaced by coke in steel production due to cost, even though coke usually adds sulphur and sometimes other deleterious contaminants to the pig iron. Wooded metallurgical regions devoid of coal like Sweden, the Urals, or Siberia transitioned from charcoal in the early 20th century.
Cooking and heating fuel
Prior to the Industrial Revolution, charcoal was occasionally used as a cooking fuel. It is counted as a smokeless fuel; that is, the carbon is sufficiently pure that burning it causes substantially less air pollution than burning the original uncarbonized organic material would. In the 20th century, clean-air legislation mandated smokeless fuels (mostly coke or charcoal) in many areas of Europe. In the 21st century, charcoal has been advocated as a way to improve the health of people burning raw biomass for cooking and/or heating. Modern "charcoal" briquettes, widely used for outdoor cooking, are made with charcoal but may also include coal as an energy source as well as accelerants, binders and filler.
To contain the charcoal and use it for cooking purposes, a barbecue grill may be used. A small Japanese charcoal grill is known as a shichirin. A brazier is a container used to burn charcoal or other solid fuel.
To start the charcoal burning is harder than starting a wood fire and charcoal lighter fluid may be employed. A chimney starter or electric charcoal starter are tools to help with starting to light charcoal.
Approximately 75% of fuel burned in Haiti is charcoal.
Reducing agent
Certain types of charcoal, such as wood charcoal, are used for reducing heated metallic oxides to their respective metals:
ZnO + C → Zn + CO
Fe2O3 + 3C → 2Fe + 3CO
Charcoal can also be used to reduce super heated steam to hydrogen (along with the formation of carbon monoxide):
C + H2O (1000 °C) → H2 + CO (Water gas)
Syngas production, automotive fuel
Like many other sources of carbon, charcoal can be used for the production of various syngas compositions; i.e., various CO + H2 + CO2 + N2 mixtures. The syngas is typically used as fuel, including automotive propulsion, or as a chemical feedstock.
In times of scarce petroleum, automobiles and even buses have been converted to burn wood gas: a gas mixture consisting primarily of diluting atmospheric nitrogen, but also containing combustible gasses (mostly carbon monoxide) released by burning charcoal or wood in a wood gas generator. In 1931, Tang Zhongming developed an automobile powered by charcoal, and these cars were popular in China until the 1950s, and in occupied France during World War II, where they were called gazogènes.
Pyrotechnics
Charcoal is used in the production of black powder, which is used extensively in the production of fireworks. It is usually ground into a fine powder, with air float grade being the finest particle size available commercially. When used in black powder compositions, it is often ball-milled with other ingredients so that they are intimately mixed together. Certain charcoals perform better when used to make black powder; these include spruce, willow, paulownia and grapevine among others. Charcoal produces fine dark orange/golden sparks. Usually, powder with a mesh size from 10 to 325 is used to obtain showers of golden sparks in pyrotechnic compositions.
Cosmetic use of bamboo charcoal
Charcoal is also incorporated in multiple cosmetic products. It can be produced from regular bamboo cut into small pieces and boiled in water to remove soluble compounds. Raw bamboo charcoal is obtained after drying and carbonization in an oven at elevated temperature. The role of charcoal in cosmetics is based on its highly effective absorbing properties at a microscopic scale.
Carbon source
Charcoal may be used as a source of carbon in chemical reactions. One example of this is the production of carbon disulphide through the reaction of sulfur vapors with hot charcoal. In that case, the wood should be charred at high temperature to reduce the residual amounts of hydrogen and oxygen that lead to side reactions.
Purification and filtration
Charcoal may be activated to increase its effectiveness as a filter. Activated charcoal readily adsorbs a wide range of organic compounds dissolved or suspended in gases and liquids. In certain industrial processes, such as the purification of sucrose from cane sugar, impurities cause an undesirable color, which can be removed with activated charcoal. It is also used to absorb odors and toxins in gaseous solutions, as in home air purifiers and some types of gas mask. The medical use of activated charcoal is mainly the absorption of poisons. Activated charcoal is available without a prescription, so it is used for a variety of health-related applications. For example, it is often used to reduce discomfort and embarrassment due to excessive gas (flatulence) in the digestive tract.
Animal charcoal or bone black is the carbonaceous residue obtained by the dry distillation of bones. It contains only about 10% carbon, the remaining being calcium and magnesium phosphates (80%) and other inorganic material originally present in the bones. It is generally manufactured from the residues obtained in the glue and gelatin industries. Its bleaching power was applied in 1812 by Derosne for clarifying sugar syrup, but its use in this direction has now greatly diminished. Today it is seldom used for this purpose due to the introduction of more active and easily managed reagents, but it is still employed to some extent in laboratory practice. The bleaching action of the charcoal in solution diminishes as it adsorbs colored contaminants, and it must be reactivated periodically by separate washing and reheating. While wood charcoal effectively removes some pigments and contaminants from solutions, bone charcoal is generally more effective as an adsorption filter due to its increased porosity and surface area.
Art
Charcoal is used for drawing, making rough sketches in painting, and is one of the possible media used for making a parsemage. It usually must be preserved by the application of a fixative. Artists generally utilize charcoal in four forms:
Vine charcoal is created by burning grape vines.
Willow charcoal is created by burning willow sticks.
Powdered charcoal is often used to "tone" or cover large sections of a drawing surface. Drawing over the toned areas darkens it further, but the artist can also lighten (or completely erase) within the toned area to create lighter tones.
Compressed charcoal is charcoal powder mixed with gum binder and compressed into sticks. The amount of binder determines the hardness of the stick. Compressed charcoal is used in charcoal pencils.
Horticulture
One additional use of charcoal was rediscovered recently for horticulture. Although American gardeners have used charcoal for a short time, research on Terra preta soils in Amazonia has discovered the widespread use of biochar by pre-Columbian natives to ameliorate unproductive soil into soil rich in carbon. The technique may find modern application, both to improve soils and as a means of carbon sequestration.
Animal husbandry
Charcoal is mixed with feed, added to litter, or used in the treatment of manure. Poultry benefits from using charcoal in this manner.
A concern that activated charcoal might be used unscrupulously to allow livestock to tolerate low quality feed contaminated with aflatoxins resulted in the Association of American Feed Control Officials banning it in 2012 from use in commercial livestock feeds.
Medicine
Charcoal in the form of charcoal biscuits was consumed in the past for gastric problems. Now it can be consumed in tablet, capsule, or powder form for digestive effects. Research regarding its effectiveness is controversial.
Charcoal has been used in combination with saccharin in research to measure mucociliary transport time.
Charcoal has also been incorporated into toothpaste formulas; however, there is no evidence to determine its safety and effectiveness.
Red colobus monkeys in Africa have been observed eating charcoal for self-medication. Because their leafy diets contain high levels of cyanide, which may lead to indigestion, they learned to consume charcoal, which absorbs the cyanide and relieves discomfort. This knowledge is transmitted from mother to infant.
Environmental sustainability
Production and utilization of charcoal, like any use of woody biomass as fuel, typically results in emissions and can contribute to deforestation.
The use of charcoal as a smelting fuel has been experiencing a resurgence in South America resulting in severe environmental, social and medical problems. Charcoal production at a sub-industrial level is one of the causes of deforestation. Charcoal production is now usually illegal and nearly always unregulated, as in Brazil, where charcoal production is a large illegal industry for making pig iron.
Massive forest destruction has been documented in areas such as Virunga National Park in the Democratic Republic of Congo, where it is considered a primary threat to the survival of the mountain gorillas. Similar threats are found in Zambia. In Malawi, illegal charcoal trade employs 92,800 workers and is the main source of heat and cooking fuel for 90 percent of the nation's population. Some experts, such as Duncan MacQueen, Principal Researcher–Forest Team, International Institute for Environment and Development (IIED), argue that while illegal charcoal production causes deforestation, a regulated charcoal industry that required replanting and sustainable use of the forests "would give their people clean efficient energy – and their energy industries a strong competitive advantage".
Recent assessments of charcoal imported to Europe have shown that many charcoal products are produced from tropical wood, often of undeclared origin. In an analysis of barbecue charcoal marketed in Germany, the World Wildlife Fund found that most products contain tropical wood. As a notable exception, reference is made to barbecue charcoal imports from Namibia, where charcoal is typically produced from surplus biomass resulting from woody plant encroachment. Charcoal trafficking in Somalia is an economic and environmental issue with significant regional-security implications.
In popular culture
The last section of the film Le Quattro Volte (2010) gives a good and long, if poetic, documentation of the traditional method of making charcoal. The Arthur Ransome children's series Swallows and Amazons (particularly the second book, Swallowdale) features carefully drawn vignettes of the lives and the techniques of charcoal burners at the start of the 20th century, in the Lake District of the UK. Antonín Dvořák's opera King and Charcoal Burner is based on a Czech legend about a king who gets lost in a forest and is rescued by a charcoal burner.
| Technology | Energy and fuel | null |
6603263 | https://en.wikipedia.org/wiki/Military%20transport%20aircraft | Military transport aircraft | A military transport aircraft, military cargo aircraft or airlifter is a military-owned transport aircraft used to support military operations by airlifting troops and military equipment. Transport aircraft are crucial to maintaining supply lines to forward bases that are difficult to reach by ground or waterborne access, and can be used for both strategic and tactical missions. They are also often used for civilian emergency relief missions by transporting humanitarian aid.
Air frames
Fixed-wing
Military transport aeroplanes are defined in terms of their range capability as strategic airlift or tactical airlift to reflect the needs of the land forces which they most often support. These roughly correspond to the commercial flight length distinctions: Eurocontrol defines short-haul routes as shorter than , long-haul routes as longer than and medium-haul between.
The military glider is an unpowered tactical air transport which has been used in some campaigns to transport troops and/or equipment to the battle front.
Rotary wing
Military transport helicopters are used in places where the use of conventional aircraft is impossible. For example, the military transport helicopter is the primary transport asset of US Marines deploying from LHDs and LHA. The landing possibilities of helicopters are almost unlimited, and where landing is impossible, for example densely packed jungle, the ability of the helicopter to hover allows troops to deploy by abseiling and roping.
Transport helicopters are operated in assault, medium and heavy classes. Air assault helicopters are usually the smallest of the transport types, and designed to move an infantry squad or section and their equipment. Helicopters in the assault role are generally armed for self-protection both in transit and for suppression of the landing zone. This armament may be in the form of door gunners, or the modification of the helicopter with stub wings and pylons to carry missiles and rocket pods. For example, the Sikorsky S-70, fitted with the ESSM (External Stores Support System), and the Hip E variant of the Mil Mi-8 can carry as much disposable armament as some dedicated attack helicopters.
Medium transport helicopters are generally capable of moving up to a platoon of infantry, or transporting towed artillery or light vehicles either internally or as underslung roles. Unlike the assault helicopter they are usually not expected to land directly in a contested landing zone, but are used to reinforce and resupply landing zones taken by the initial assault wave. Examples include the unarmed versions of the Mil Mi-8, Super Puma, CH-46 Sea Knight, and NH90.
Heavy lift helicopters are the largest and most capable of the transport types, currently limited in service to the CH-53 Sea Stallion and related CH-53E Super Stallion, CH-47 Chinook, Mil Mi-26, and Aérospatiale Super Frelon. Capable of lifting up to 80 troops and moving small Armoured fighting vehicles (usually as slung loads but also internally), these helicopters operate in the tactical transport role in much the same way as small fixed wing turboprop air-lifters. The lower speed, range and increased fuel consumption of helicopters are offset by their not requiring a runway.
Payload comparison
| Technology | Military aviation | null |
6609562 | https://en.wikipedia.org/wiki/Leicester%20Longwool | Leicester Longwool | The Leicester Longwool is an English breed of sheep. Alternative names for the breed include: Leicester, Bakewell Leicester, Dishley Leicester, English Leicester, Improved Leicester and New Leicester. It was originally developed by 18th-century breeding innovator Robert Bakewell.
It is now one of Britain's rarest breeds, categorised as "endangered" by the Rare Breeds Survival Trust, since fewer than 500 registered breeding females remain in the United Kingdom.
History
Leicester Longwool sheep date back to the 1700s, and were found in the Midland counties of England, originally developed in Dishley Grange, Leicestershire, by Robert Bakewell. Bakewell was the foremost exponent of modern animal-breeding techniques in the selection of livestock.
The Leicester Longwool in the 1700s was slow-growing and coarsely boned. They now have been developed to gain weight quickly and are fast-growing.
Leicester Longwool was one of the first pure sheep breeds introduced to Australia, having been introduced in 1826.
The Leicester Longwool has been used to improve many sheep breeds because of its meaty carcase (carcass) and heavy fleece.
Characteristics
The head of an Leicester Longwool should be carried well, not too high and should have no signs of horns on the poll (forehead). The face is generally in a wedge shape, covered in white hairs and can appear to have a blue tinge. The lips and nostrils should be black. Having black specks on the face and ears is not objectionable. The neck should be of medium length. The shoulders should be strong and level with the back, which should be flat. The legs should be straight and wide apart and the hooves should be black.
The fleece should be dense (having thick and blocky clumps of wool also known as the staple). It should be lustrous, indicating the shine on the wool, and should have a well-defined crimp or wave from skin to tip. The common fibre diameter for an Leicester Longwool is 32 to 38 micrometres (microns).
The Leicester Longwool should be free, active and well balanced while in movement. It should appear to be alert and robust, showing style and character.
Other information
Leicester Longwool sheep are currently found in Australia, New Zealand, Great Britain, Sweden and the United States. They are sound-footed, which means they are acceptable on flat, hilly or more especially marginal country.
They are large-framed with wide, even toplines (backs), strong constitutions and good temperaments.
| Biology and health sciences | Sheep | Animals |
6613070 | https://en.wikipedia.org/wiki/System%20time | System time | In computer science and computer programming, system time represents a computer system's notion of the passage of time. In this sense, time also includes the passing of days on the calendar.
System time is measured by a system clock, which is typically implemented as a simple count of the number of ticks that have transpired since some arbitrary starting date, called the epoch. For example, Unix and POSIX-compliant systems encode system time ("Unix time") as the number of seconds elapsed since the start of the Unix epoch at 1 January 1970 00:00:00 UT, with exceptions for leap seconds. Systems that implement the 32-bit and 64-bit versions of the Windows API, such as Windows 9x and Windows NT, provide the system time as both , represented as a year/month/day/hour/minute/second/milliseconds value, and , represented as a count of the number of 100-nanosecond ticks since 1 January 1601 00:00:00 UT as reckoned in the proleptic Gregorian calendar.
System time can be converted into calendar time, which is a form more suitable for human comprehension. For example, the Unix system time seconds since the beginning of the epoch translates into the calendar time 9 September 2001 01:46:40 UT. Library subroutines that handle such conversions may also deal with adjustments for time zones, daylight saving time (DST), leap seconds, and the user's locale settings. Library routines are also generally provided that convert calendar times into system times.
Many implementations that currently store system times as 32-bit integer values will suffer from the impending Year 2038 problem. These time values will overflow ("run out of bits") after the end of their system time epoch, leading to software and hardware errors. These systems will require some form of remediation, similar to efforts required to solve the earlier Year 2000 problem. This will also be a potentially much larger problem for existing data file formats that contain system timestamps stored as 32-bit values.
Other time measurements
Closely related to system time is process time, which is a count of the total CPU time consumed by an executing process. It may be split into user and system CPU time, representing the time spent executing user code and system kernel code, respectively. Process times are a tally of CPU instructions or clock cycles and generally have no direct correlation to wall time.
File systems keep track of the times that files are created, modified, and/or accessed by storing timestamps in the file control block (or inode) of each file and directory.
History
Most first-generation personal computers did not keep track of dates and times. These included systems that ran the CP/M operating system, as well as early models of the Apple II, the BBC Micro, and the Commodore PET, among others. Add-on peripheral boards that included real-time clock chips with on-board battery back-up were available for the IBM PC and XT, but the IBM AT was the first widely available PC that came equipped with date/time hardware built into the motherboard. Prior to the widespread availability of computer networks, most personal computer systems that did track system time did so only with respect to local time and did not make allowances for different time zones.
With current technology, most modern computers keep track of local civil time, as do many other household and personal devices such as VCRs, DVRs, cable TV receivers, PDAs, pagers, cell phones, fax machines, telephone answering machines, cameras, camcorders, central air conditioners, and microwave ovens.
Microcontrollers operating within embedded systems (such as the Raspberry Pi, Arduino, and other similar systems) do not always have internal hardware to keep track of time. Many such controller systems operate without knowledge of the external time. Those that require such information typically initialize their base time upon rebooting by obtaining the current time from an external source, such as from a time server or external clock, or by prompting the user to manually enter the current time.
Implementation
The system clock is typically implemented as a programmable interval timer that periodically interrupts the CPU, which then starts executing a timer interrupt service routine. This routine typically adds one tick to the system clock (a simple counter) and handles other periodic housekeeping tasks (preemption, etc.) before returning to the task the CPU was executing before the interruption.
Retrieving system time
The following tables illustrate methods for retrieving the system time in various operating systems, programming languages, and applications. Values marked by (*) are system-dependent and may differ across implementations. All dates are given as Gregorian or proleptic Gregorian calendar dates.
The resolution of an implementation's measurement of time does not imply the same precision of such measurements. For example, a system might return the current time as a value measured in microseconds, but actually be capable of discerning individual clock ticks with a frequency of only 100 Hz (10 ms).
Operating systems
Programming languages and applications
| Technology | Computer architecture concepts | null |
6614349 | https://en.wikipedia.org/wiki/Plant%20reproduction | Plant reproduction | Plant reproduction is the production of new offspring in plants, which can be accomplished by sexual or asexual reproduction. Sexual reproduction produces offspring by the fusion of gametes, resulting in offspring genetically different from either parent. Asexual reproduction produces new individuals without the fusion of gametes, resulting in clonal plants that are genetically identical to the parent plant and each other, unless mutations occur.
Asexual reproduction
Asexual reproduction does not involve the production and fusion of male and female gametes. Asexual reproduction may occur through budding, fragmentation, spore formation, regeneration and vegetative propagation.
Asexual reproduction is a type of reproduction where the offspring comes from one parent only, thus inheriting the characteristics of the parent. Asexual reproduction in plants occurs in two fundamental forms, vegetative reproduction and agamospermy. Vegetative reproduction involves a vegetative piece of the original plant producing new individuals by budding, tillering, etc. and is distinguished from apomixis, which is a replacement of sexual reproduction, and in some cases involves seeds. Apomixis occurs in many plant species such as dandelions (Taraxacum species) and also in some non-plant organisms. For apomixis and similar processes in non-plant organisms, see parthenogenesis.
Natural vegetative reproduction is a process mostly found in perennial plants, and typically involves structural modifications of the stem or roots and in a few species leaves. Most plant species that employ vegetative reproduction do so as a means to perennialize the plants, allowing them to survive from one season to the next and often facilitating their expansion in size. A plant that persists in a location through vegetative reproduction of individuals gives rise to a clonal colony. A single ramet, or apparent individual, of a clonal colony is genetically identical to all others in the same colony. The distance that a plant can move during vegetative reproduction is limited, though some plants can produce ramets from branching rhizomes or stolons that cover a wide area, often in only a few growing seasons. In a sense, this process is not one of reproduction but one of survival and expansion of biomass of the individual. When an individual organism increases in size via cell multiplication and remains intact, the process is called vegetative growth. However, in vegetative reproduction, the new plants that result are new individuals in almost every respect except genetic. A major disadvantage of vegetative reproduction is the transmission of pathogens from parent to offspring. It is uncommon for pathogens to be transmitted from the plant to its seeds (in sexual reproduction or in apomixis), though there are occasions when it occurs.
Seeds generated by apomixis are a means of asexual reproduction, involving the formation and dispersal of seeds that do not originate from the fertilization of the embryos. Hawkweeds (Hieracium), dandelions (Taraxacum), some species of Citrus and Kentucky blue grass (Poa pratensis) all use this form of asexual reproduction. Pseudogamy occurs in some plants that have apomictic seeds, where pollination is often needed to initiate embryo growth, though the pollen contributes no genetic material to the developing offspring. Other forms of apomixis occur in plants also, including the generation of a plantlet in replacement of a seed or the generation of bulbils instead of flowers, where new cloned individuals are produced.
Structures
A rhizome is a modified underground stem serving as an organ of vegetative reproduction; the growing tips of the rhizome can separate as new plants, e.g., polypody, iris, couch grass and nettles.
Prostrate aerial stems, called runners or stolons, are important vegetative reproduction organs in some species, such as the strawberry, numerous grasses, and some ferns.
Adventitious buds form on roots near the ground surface, on damaged stems (as on the stumps of cut trees), or on old roots. These develop into above-ground stems and leaves. A form of budding called suckering is the reproduction or regeneration of a plant by shoots that arise from an existing root system. Species that characteristically produce suckers include elm (Ulmus) and many members of the rose family such as Rosa, Kerria and Rubus.
Bulbous plants such as onion (Allium cepa), hyacinths, narcissi and tulips reproduce vegetatively by dividing their underground bulbs into more bulbs. Other plants like potatoes (Solanum tuberosum) and dahlias reproduce vegetatively from underground tubers. Gladioli and crocuses reproduce vegetatively in a similar way with corms.
Gemmae are single cells or masses of cells that detach from plants to form new clonal individuals. These are common in Liverworts and mosses and in the gametophyte generation of some filmy fern. They are also present in some Club mosses such as Huperzia lucidula . They are also found in some higher plants such as species of Drosera.
Usage
The most common form of plant reproduction used by people is seeds, but a number of asexual methods are used which are usually enhancements of natural processes, including: cutting, grafting, budding, layering, division, sectioning of rhizomes, roots, tubers, bulbs, stolons, tillers, etc., and artificial propagation by laboratory tissue cloning. Asexual methods are most often used to propagate cultivars with individual desirable characteristics that do not come true from seed. Fruit tree propagation is frequently performed by budding or grafting desirable cultivars (clones), onto rootstocks that are also clones, propagated by stooling.
In horticulture, a cutting is a branch that has been cut off from a mother plant below an internode and then rooted, often with the help of a rooting liquid or powder containing hormones. When a full root has formed and leaves begin to sprout anew, the clone is a self-sufficient plant, genetically identical.
Examples include cuttings from the stems of blackberries (Rubus occidentalis), African violets (Saintpaulia), verbenas (Verbena) to produce new plants. A related use of cuttings is grafting, where a stem or bud is joined onto a different stem. Nurseries offer for sale trees with grafted stems that can produce four or more varieties of related fruits, including apples. The most common usage of grafting is the propagation of cultivars onto already rooted plants, sometimes the rootstock is used to dwarf the plants or protect them from root damaging pathogens.
Since vegetatively propagated plants are clones, they are important tools in plant research. When a clone is grown in various conditions, differences in growth can be ascribed to environmental effects instead of genetic differences.
Sexual reproduction
Sexual reproduction involves two fundamental processes: meiosis, which rearranges the genes and reduces the number of chromosomes, and fertilization, which restores the chromosome to a complete diploid number. In between these two processes, different types of plants and algae vary, but many of them, including all land plants, undergo alternation of generations, with two different multicellular structures (phases), a gametophyte and a sporophyte. The evolutionary origin and adaptive significance of sexual reproduction are discussed in the pages Evolution of sexual reproduction and Origin and function of meiosis.
The gametophyte is the multicellular structure (plant) that is haploid, containing a single set of chromosomes in each cell. The gametophyte produces male or female gametes (or both), by a process of cell division, called mitosis. In vascular plants with separate gametophytes, female gametophytes are known as mega gametophytes (mega=large, they produce the large egg cells) and the male gametophytes are called micro gametophytes (micro=small, they produce the small sperm cells).
The fusion of male and female gametes (fertilization) produces a diploid zygote, which develops by mitotic cell divisions into a multicellular sporophyte.
The mature sporophyte produces spores by meiosis, sometimes referred to as reduction division because the chromosome pairs are separated once again to form single sets.
In mosses and liverworts, the gametophyte is relatively large, and the sporophyte is a much smaller structure that is never separated from the gametophyte. In ferns, gymnosperms, and flowering plants (angiosperms), the gametophytes are relatively small and the sporophyte is much larger. In gymnosperms and flowering plants the megagametophyte is contained within the ovule (that may develop into a seed) and the microgametophyte is contained within a pollen grain.
History of sexual reproduction of plants
Unlike animals, plants are immobile, and cannot seek out sexual partners for reproduction. In the evolution of early plants, abiotic means, including water and much later, wind, transported sperm for reproduction. The first plants were aquatic, as described in the page Evolutionary history of plants, and released sperm freely into the water to be carried with the currents. Primitive land plants such as liverworts and mosses had motile sperm that swam in a thin film of water or were splashed in water droplets from the male reproduction organs onto the female organs. As taller and more complex plants evolved, modifications in the alternation of generations evolved. In the Paleozoic era progymnosperms reproduced by using spores dispersed on the wind. The seed plants including seed ferns, conifers and cordaites, which were all gymnosperms, evolved 350 million years ago. They had pollen grains that contained the male gametes for protection of the sperm during the process of transfer from the male to female parts.
It is believed that insects fed on the pollen, and plants thus evolved to use insects to actively carry pollen from one plant to the next. Seed producing plants, which include the angiosperms and the gymnosperms, have a heteromorphic alternation of generations with large sporophytes containing much-reduced gametophytes. Angiosperms have distinctive reproductive organs called flowers, with carpels, and the female gametophyte is greatly reduced to a female embryo sac, with as few as eight cells. Each pollen grains contains a greatly reduced male gametophyte consisting of three or four cells. The sperm of seed plants are non-motile, except for two older groups of plants, the Cycadophyta and the Ginkgophyta, which have flagella.
Flowering plants
Flowering plants, the dominant plant group, reproduce both by sexual and asexual means. Their distinguishing feature is that their reproductive organs are contained in flowers. Sexual reproduction in flowering plants involves the production of separate male and female gametophytes that produce gametes.
The anther produces pollen grains that contain male gametophytes. The pollen grains attach to the stigma on top of a carpel, in which the female gametophytes (inside ovules) are located. Plants may either self-pollinate or cross-pollinate. The transfer of pollen (the male gametophytes) to the female stigmas occurs is called pollination. After pollination occurs, the pollen grain germinates to form a pollen tube that grows through the carpel's style and transports male nuclei to the ovule to fertilize the egg cell and central cell within the female gametophyte in a process termed double fertilization. The resulting zygote develops into an embryo, while the triploid endosperm (one sperm cell plus a binucleate female cell) and female tissues of the ovule give rise to the surrounding tissues in the developing seed. The fertilized ovules develop into seeds within a fruit formed from the ovary. When the seeds are ripe they may be dispersed together with the fruit or freed from it by various means to germinate and grow into the next generation.
Pollination
Plants that use insects or other animals to move pollen from one flower to the next have developed greatly modified flower parts to attract pollinators and to facilitate the movement of pollen from one flower to the insect and from the insect to the next flower. Flowers of wind-pollinated plants tend to lack petals and or sepals; typically large amounts of pollen are produced and pollination often occurs early in the growing season before leaves can interfere with the dispersal of the pollen. Many trees and all grasses and sedges are wind-pollinated.
Plants have a number of different means to attract pollinators including color, scent, heat, nectar glands, edible pollen and flower shape. Along with modifications involving the above structures two other conditions play a very important role in the sexual reproduction of flowering plants, the first is the timing of flowering and the other is the size or number of flowers produced. Often plant species have a few large, very showy flowers while others produce many small flowers, often flowers are collected together into large inflorescences to maximize their visual effect, becoming more noticeable to passing pollinators. Flowers are attraction strategies and sexual expressions are functional strategies used to produce the next generation of plants, with pollinators and plants having co-evolved, often to some extraordinary degrees, very often rendering mutual benefit.
The largest family of flowering plants is the orchids (Orchidaceae), estimated by some specialists to include up to 35,000 species, which often have highly specialized flowers that attract particular insects for pollination. The stamens are modified to produce pollen in clusters called pollinia, which become attached to insects that crawl into the flower. The flower shapes may force insects to pass by the pollen, which is "glued" to the insect. Some orchids are even more highly specialized, with flower shapes that mimic the shape of insects to attract them to attempt to 'mate' with the flowers, a few even have scents that mimic insect pheromones.
Another large group of flowering plants is the Asteraceae or sunflower family with close to 22,000 species, which also have highly modified inflorescences composed of many individual flowers called florets. Heads with florets of one sex, when the flowers are pistillate or functionally staminate or made up of all bisexual florets, are called homogamous and can include discoid and liguliflorous type heads. Some radiate heads may be homogamous too. Plants with heads that have florets of two or more sexual forms are called heterogamous and include radiate and disciform head forms.
Ferns
Ferns typically produce large diploids with stem, roots, and leaves. On fertile leaves sporangia are produced, grouped together in sori and often protected by an indusium. If the spores are deposited onto a suitable moist substrate they germinate to produce short, thin, free-living gametophytes called prothalli that are typically heart-shaped, small and green in color. The gametophytes produce both motile sperm in the antheridia and egg cells in separate archegonia. After rains or when dew deposits a film of water, the motile sperm are splashed away from the antheridia, which are normally produced on the top side of the thallus, and swim in the film of water to the antheridia where they fertilize the egg. To promote out crossing or cross-fertilization the sperm is released before the eggs are receptive of the sperm, making it more likely that the sperm will fertilize the eggs of the different thallus. A zygote is formed after fertilization, which grows into a new sporophytic plant. The condition of having separate sporophyte and gametophyte plants is called alternation of generations. Other plants with similar reproductive strategies include Psilotum, Lycopodium, Selaginella and Equisetum.
Bryophytes
The bryophytes, which include liverworts, hornworts and mosses, can reproduce both sexually and vegetatively. The life cycles of these plants start with haploid spores that grow into the dominant form, which is a multicellular haploid gametophyte, with thalloid or leaf-like structures that photosynthesize. The gametophyte is the most commonly known phase of the plant. Bryophytes are typically small plants that grow in moist locations and like ferns, have motile sperm which swim to the ovule using flagella and therefore need water to facilitate sexual reproduction. Bryophytes show considerable variation in their reproductive structures, and a basic outline is as follows: Haploid gametes are produced in antheridia and archegonia by mitosis. The sperm released from the antheridia respond to chemicals released by ripe archegonia and swim to them in a film of water and fertilize the egg cells, thus producing zygotes that are diploid. The zygote divides repeatedly by mitotic division and grows into a diploid sporophyte. The resulting multicellular diploid sporophyte produces spore capsules called sporangia. The spores are produced by meiosis, and when ripe, the capsules burst open to release the spores. In some species each gametophyte is one sex while other species may be monoicous, producing both antheridia and archegonia on the same gametophyte which is thus hermaphrodite.
Dispersal and offspring care
One of the outcomes of plant reproduction is the generation of seeds, spores, and fruits that allow plants to move to new locations or new habitats.
Plants do not have nervous systems or any will for their actions. Even so, scientists are able to observe mechanisms that help their offspring thrive as they grow. All organisms have mechanisms to increase survival in offspring.
Offspring care is observed in the Mammillaria hernandezii, a small cactus found in Mexico. A cactus is a type of succulent, meaning it retains water when it is available for future droughts. M. hernandezii also stores a portion of its seeds in its stem, and releases the rest to grow. This can be advantageous for many reasons. By delaying the release of some of its seeds, the cactus can protect these from potential threats from insects, herbivores, or mold caused by micro-organisms. A study found that the presence of adequate water in the environment causes M. Hernandezii to release more seeds to allow for germination. The plant was able to perceive a water potential gradient in the surroundings, and act by giving its seeds a better chance in this preferable environment. This evolutionary strategy gives a better potential outcome for seed germination.
| Biology and health sciences | Plant reproduction | null |
3708293 | https://en.wikipedia.org/wiki/Vechur%20cattle | Vechur cattle | Vechur is a breed of zebu (Bos indicus) cattle, named after the village of Vechoor in Kerala, India. With an average length of 124 cm and height of 87 cm, it is the smallest cattle breed in the world according to the Guinness Book of Records, and is valued for the larger amount of milk it produces relative to the amount of food it requires.
History
The Vechur animals were saved from extinction due to conservation efforts by Sosamma Iype, a Professor of Animal breeding and Genetics along with a team of her students. In 1989, a conservation unit was started. A Conservation trust was formed in 1998 to continue the work with farmer participation.
The Vechur cow was popular in Kerala until the 1960s, but became rare when native cattle were crossbred with exotic varieties.
In 2000, the Vechur cow was listed on the FAO's World Watch List of Domestic Animal Diversity, in its ‘Critical-Maintained Breeds List', pointing to imminent extinction as breeds are included in the list when the number of breeding females and males fall to very low levels. About 200 cows are supposed to exist today, nearly 100 of them with the Veterinary College.
Characteristics
The breed averages about 90 cm in height and weighs around 130 kg, yielding up to 3 litres of milk a day. This is much less than that of the hybrid varieties but, unlike them, the Vechur cow also requires little by way of feed or maintenance. The milk is believed to have medicinal qualities and easy digestibility due to smaller fat globule size. A recent report claims that the milk of the Vechur cow has more of the beta casein variety A2, rather than the variety A1 which is implicated in conditions like diabetes, ischaemic heart disease and autism.
The medicinal properties of the Vechur cow's milk have been documented traditionally in Ayurveda and recent scientific studies have substantiated this. The protein component of the Vechur cow's milk has an improved antimicrobial property. According to recent findings, the anti-bacterial property of the lactoferrin protein present in the Vechur Cow milk is more than that of the antibiotic ampicillin. Although lactoferrin, present in the milk of all mammals, is found to have antimicrobial, antiviral, antitumor, immunodefence and anti-inflammatory properties recent studies prove that in the case of the lactoferrin protein in the Vechur Cow milk, these properties are much more enhanced. The Vechur Ghee (clarified butter) produced from Vechur cow’s milk, is famous for its high medicinal values due to the presence of A2 beta-lactalbumin protein and higher arginine content which is good for the health of convalescing people.
Controversies
A controversy arose in 1997 when environmentalist Vandana Shiva alleged that a Scottish company, The Roslin Institute was trying to patent the cow's genetic code. Shiva described the action as piracy. The Roslin Institute denied the charge which was subsequently proved to be baseless.
| Biology and health sciences | Miniature cattle | Animals |
3709072 | https://en.wikipedia.org/wiki/Gregorian%20telescope | Gregorian telescope | The Gregorian telescope is a type of reflecting telescope designed by Scottish mathematician and astronomer James Gregory in the 17th century, and first built in 1673 by Robert Hooke. James Gregory was a contemporary of Isaac Newton, and both often worked simultaneously on similar projects. Gregory's design was published in 1663 and pre-dates the first practical reflecting telescope, the Newtonian telescope, built by Sir Isaac Newton in 1668. However, Gregory's design was only a theoretical description, and he never actually constructed the telescope. It was not successfully built until five years after Newton's first reflecting telescope.
History
The Gregorian telescope is named after the James Gregory design, which appeared in his 1663 publication (The Advance of Optics). Similar theoretical designs have been found in the writings of Bonaventura Cavalieri ( (On Burning Mirrors), 1632) and Marin Mersenne (, 1636). Gregory's early attempts to build the telescope failed, since he had no practical skill himself and could find no optician capable of actually constructing one. It was not until ten years after Gregory's publication, aided by the interest of experimental scientist Robert Hooke, that a working instrument was created. The early Scottish optician and telescope maker James Short built Gregorian telescopes with parabolic mirrors made from the highly reflective speculum metal.
Design
The Gregorian telescope consists of two concave mirrors: the primary mirror (a concave paraboloid) collects the light and brings it to a focus before the secondary mirror (a concave ellipsoid), where it is reflected back through a hole in the centre of the primary, and thence out the bottom end of the instrument, where it can be viewed with the aid of the eyepiece.
The Gregorian design solved the problem of viewing the image in a reflector by allowing the observer to stand behind the primary mirror. This design of telescope renders an erect image, making it useful for terrestrial observations. It also works as a telephoto lens with its tube much shorter than the system's actual focal length.
The design was largely superseded by the Cassegrain telescope. It is still used for some spotting scopes because this design creates an erect image without the need for prisms. The Steward Observatory Mirror Lab has been making mirrors for large Gregorian telescopes at least since 1985.
In the Gregorian design, the primary mirror creates a real image before the secondary mirror. This allows for a field stop to be placed at this location, so that the light from outside the field of view does not reach the secondary mirror. This is a major advantage for solar telescopes, where a field stop (Gregorian stop) can reduce the amount of heat reaching the secondary mirror and subsequent optical components. The Solar Optical Telescope on the Hinode satellite is one example of this design.
For amateur telescope makers the Gregorian can be less difficult to fabricate than a Cassegrain because the concave secondary is Foucault-testable like the primary, which is not the case with the Cassegrain's convex secondary.
Gallery
Examples
The MeerKAT, the Green Bank Telescope, the Arecibo Observatory, and the Allen Telescope Array are all radio telescopes employing off-axis Gregorian optics.
The Vatican Advanced Technology Telescope, the Magellan telescopes, and the Large Binocular Telescope use Gregorian optics.
The Giant Magellan Telescope will also use Gregorian optics.
The NSF's Daniel K. Inouye Solar Telescope
| Technology | Telescope | null |
3716072 | https://en.wikipedia.org/wiki/Bull%20and%20terrier | Bull and terrier | Bull and terrier was a common name for crossbreeds between bulldogs and terriers in the early 1800s. Other names included half-and-halfs and half-breds. It was a time in history when, for thousands of years, dogs were classified by use or function, unlike the modern pets of today that were bred to be conformation show dogs and family pets. Bull and terrier crosses were originally bred to function as fighting dogs for bull- and bear-baiting, and other popular blood sports during the Victorian era. The sport of bull baiting required a dog with attributes such as tenacity and courage, a wide frame with heavy bone, and a muscular, protruding jaw. By crossing bulldogs with various terriers from Ireland and Great Britain, breeders introduced "gameness and agility" into the hybrid mix.
Little is known about the pedigrees of bull and terrier crosses, or any other crosses that originated during that time. The types and styles of dogs varied geographically depending on individual preferences. Breeders in one area may have preferred a cross with a higher percentage of terrier than bulldog. Some early anecdotal reports indicate that bulldog to terrier was preferred over bull and terrier to bull terrier, which was likely to have resulted in at least half or more bulldog blood. The bull and terrier was never a bona fide breed; rather, it referred to a heterogeneous group of dogs that may include purebreds involving different breeds, as well as dogs believed to be crosses of those breeds. Those crossbreeds or hybrids are considered the forerunner of several modern standardised breeds.
In the mid-1830s, when enforcement of the ban on bull baiting had begun, the popularity of the original purebred bulldogs declined, and a major shift in canine genetics was occurring. The appearance of certain dogs were being altered by crossbreeding to better suit function. Not only were appearances of dogs changing, so was the terminology used to describe various breeds and dog types as recorded in ancient records. Such changes began casting doubts over the bulldog's earliest known ancestors.
Terminology
One example of how changing terminology over the centuries has caused confusion is the ubiquitous misuse of descriptors. For example, mastiff is a common descriptor for all large dogs, which created a cloud over the earliest origins of the bulldog. The Alaunt was once believed to be the likely ancestor of bulldogs and mastiffs, both of which came from Asia; others believed bulldogs descended only from mastiffs.
Over the centuries, hybrid bull and terrier crosses have been labelled with several aliases, such as half-and-half, half-bred, pit dog, bulldog terrier and pit bulldog. The most popular name was bull-terrier, a name later applied to the breed James Hinks was developing in the latter half of the 19th century. There are also many paintings, texts, and engravings created during or prior to this period that labelled the bull-and-terrier only as "bull-terrier". Hinks was still developing his new bull terrier, nicknamed White Cavalier, which he presented at the Birmingham show in May 1862.
The term pit bull terrier was sometimes applied, though later applied when naming the American Pit Bull Terrier, a modern standardised breed. The term "pit bull" is a ubiquitous term that is often misused to infer that the pit bull is a bona fide breed of dog, when it actually refers to a diverse group of dogs that may include purebred dogs of many breeds as well as dogs that are assumed to be blends of those breeds. These types of descriptors vary, depending on the recognised breeds and observers' perspectives. Despite anecdotal misinformation and incorrect visual identification, dog owners, animal shelters, veterinarians and the general public routinely use the term "pit bull" in casual and official papers as though it denotes a single, recognised breed.
History
Bulldogs of the 1800s were described as having a "round head, short nose, small ears and wide, muscular frame and legs." Their temperament has been described as being more aggressive and ferocious than other dogs of the time. It is believed their ancestors may have been mastiffs of Asian descent because of their aggressive tendencies and strength, but the term "mastiff" was used as a general reference for large dogs. Whether or not they were crossed with Pugs to make them better at bull baiting remains controversial.
It is assumed that bull and terrier hybrids were crossed with several varieties of bulldogs and terriers, the types of which depended on location. A 2016 genetic assessment verified that bulldogs were descendants of mastiffs, but it also discovered pugs in the cross. The assessment, which analysed a particular group of individual English bulldogs, used DNA rather than pedigrees to confirm that genetic diversity actually still exists. It further confirmed a substantial loss of genetic diversity in the breed resulting from a small founder population of about 68 individuals. The impact of focused selection for breeding dogs with specific physical traits created artificial genetic bottlenecks.
In Ireland, they used the old Irish bulldog with different terriers and some insertion of hunting sighthound/terrier crosses. In England, there were several varieties such as the Walsall and the Cradley Heath types. Phil Drabble reported that among the various types of bull and terrier, the type from Cradley Heath was recognised as a separate breed to be named the Staffordshire Bull Terrier. In the 19th century, the Walsall type was carried by immigrants to the United States, where it served as an important component for the genetic basis of the American Pit Bull Terrier breed, through specimens such as the dog Lloyd's Pilot and the Colby bloodline, strongly combined with Irish strains. The anatomy of the bull and terrier is the result of selective breeding for the purpose of hunting, dog fighting and baiting.
Descendants
In "Popular and Illustrated Dog Encyclopaedia" (1934–35), Major Mitford Price wrote, "The original Bull Terrier, or Bull-and-Terrier, as he was then styled, bred for fighting in the pits, bore a far closer resemblance to the Bulldog of that day than to his terrier forebears: for there exists scores of old prints as evidence that the old Bulldog, as well as the Bull-and-Terrier had the unexaggerated (in comparison with the absurd modern standards) Bulldog head, and the legs, straight and longer, of the terrier. At the same time that the new Bull-and-Terrier made its appearance, the Bulldog fanciers began breeding their animals heavier and lower to the ground, so that the Bulldog acquired a new type.
Six distinct breeds descended from the bull and terrier hybrids, five of which were recognised by the American Kennel Club (AKC) in the following order: Bull Terrier, Boston Terrier, American Staffordshire Terrier (AmStaff), Staffordshire Bull Terrier, Miniature Bull Terrier. All five breeds have also been recognised by the Canadian Kennel Club (CKC), and Fédération cynologique internationale (FCI). The American Pit Bull Terrier (APBT) is recognised by the United Kennel Club (UKC), and American Dog Breeders Association. The AKC does not recognise the American Pit Bull Terrier. After being petitioned in the 1930's to recognise the breed, they relented, however; they recognised these dogs under a different name, the Staffordshire Terrier. The name was later changed to American Staffordshire Terrier, to prevent the dogs from being confused with the English Staffordshire Terriers. The American Staffordshire Terrier and the American Pit Bull Terrier were once the same breed, and many still consider them to be. Today, there are some dogs which have been dual registered as both an APBT through the UKC, and as an AmStaff through the AKC.
DNA analysis
Geneticists have been able to further refine the sparse historical aspects of breed formation, and the time of hybridisation. A 2017 genome-wide research study suggests the following: "In this analysis, all of the bull and terrier crosses map to the terriers of Ireland and date to 1860-1870. This coincides perfectly with the historical descriptions that, though they do not clearly identify all breeds involved, report the popularity of dog contests in Ireland and the lack of stud book veracity, hence undocumented crosses, during this era of breed creation (Lee, 1894)."
It also confirms that the bull and terrier was a heterogeneous group of dogs that may include purebreds involving different breeds, as well as dogs believed to be crosses of those breeds. By 1874, in Britain the first Kennel Club Stud Book was published, which included Bull Terriers and Bulldogs.
Hunters
Some believe that the courage of most terriers, both past and present, to bear the bites of badgers and other prey they are meant to corner, dig for, or attack is derived from having a quarter to an eighth of Old English Bulldog ancestry. Other terriers that were not developed by crossing Old English Bulldogs and earth-working dogs were believed to be of inferior quality and were valued far less. There are earth-working dogs who by default and definition are called terriers because they have the ability to go to ground; however, the best earth-working and hunting terriers were regarded as the progeny of bulldogs bred to earth-working dogs (terriers), with the offspring known as bull-terriers or half-bred dogs, among other names.
Walsh also wrote of the Fox Terrier (or, rather, of its cross-breed ancestor): "The field fox-terrier, used for bolting the fox when gone to ground, was of this breed [bull and terrier]." Not only is the Fox Terrier the progeny of the bull-and-terrier, but so is the Airedale Terrier, rat-working terriers, working black and tan terriers and most all other vermin-hunting terriers. James Rodwell described in his book titled The Rat: Its History and Destructive Nature, that the great object, among the various breeders of bull-and-terrier dogs for hunting vermin and rats, was to have them as nearly thorough-bred bull as possible, but at the same time preserving all the outward appearances of the terrier as to size, shape and colour.
Dog fighting
[[File:A Dog Fight at Kit Burn’s.jpg|thumb|230px|A Dog Fight at Kit Burns''' by Edward Winslow Martin. New York, 1868.]]
In the 19th century, breeders crossed Bulldogs and English White Terriers to produce fighting dogs that were the forebears of the modern Bull Terrier breed. When blood sports were banned in the early 1800s, breeders continued with their clandestine dog fights in discreet venues, such as basements and warehouses. A major shift in canine genetics occurred during the Victorian era, at which time the appearance of certain dogs were being actively altered. The early bulldogs of the 1800s were described as having a "round head, short nose, small ears and wide, muscular frame and legs." In the 1830s, the ban on bull baiting was being strictly enforced, and with it, a noticeable decline in the popularity of the original bulldogs. Breeders had already begun crossing bulldogs with terriers for clandestine pit fighting.
James Hinks, a dog breeder in Birmingham, is credited for his role in helping to standardise the bull and terrier hybrid. Hinks introduced new blood, presumably Collies to add length to the muzzle. His version was becoming a more streamlined version of the bulldog and terrier hybrid while still maintaining substance. Hink's son said that, early on, his father also used Dalmatians to create the Bull Terrier's striking all-white coat. Others have suggested that Hinks straightened the bulldog's tendency for bowed legs by adding Pointer blood, or possibly Greyhound. Hink's son recalled, "In short, they became the old fighting dog civilized, with all of his rough edges smoothed down without being softened; alert, active, plucky, muscular, and a real gentleman." Hink's early Bull Terriers were white which gave rise to congenital sensorineural deafness (CSD), a genetic condition linked to coat colour phenotypes in English bull terriers with genetic variations that go beyond coat colour. The appearance of the Bull Terrier continued to change over time, and by the 20th century its egg-shaped head had become more prominent, soon to be standardised along with the various colours that had been introduced.
Famous bull and terriers
Author David Harris describes in his book The Bully Breeds, two illustrations of prize fighting dogs. The first was Trusty appearing in an 1806 issue of Sporting Magazine, and the second was Dustman appearing in an 1812 issue of the same magazine. Trusty was a fawn coloured bull and terrier that showed more bulldog traits than terrier, and was reputed to be undefeated, having won 104 dog fights. Trusty was purchased by Thomas Pitt, 2nd Baron Camelford and presented to Jem Belcher, a champion prize fighter of England. Unlike Trusty, Dustman appeared to have more terrier traits than bulldog and was used for badger-baiting.
In 1822, Pierce Egan, a sporting event commentator of the 1820s, first introduced the name Bull Terrier (not to be confused to the 1880s Hinks' Bull Terrier). Subsequently, Bulls-eye was introduced in The Adventures of Oliver Twist (1838), presumed to be a bull terrier owned by the villain Bill Sykes.
Another notable bull and terrier was named Billy'', weighing approximately 26 pounds; "his most notable feat was killing 100 rats in 5 min 30 sec at the Cockpit in Tufton Street, Westminster, London, UK on 23 April 1825."
| Biology and health sciences | Dogs | Animals |
20114039 | https://en.wikipedia.org/wiki/Langmuir%20adsorption%20model | Langmuir adsorption model | The Langmuir adsorption model explains adsorption by assuming an adsorbate behaves as an ideal gas at isothermal conditions. According to the model, adsorption and desorption are reversible processes. This model even explains the effect of pressure; i.e., at these conditions the adsorbate's partial pressure is related to its volume adsorbed onto a solid adsorbent. The adsorbent, as indicated in the figure, is assumed to be an ideal solid surface composed of a series of distinct sites capable of binding the adsorbate. The adsorbate binding is treated as a chemical reaction between the adsorbate gaseous molecule and an empty sorption site . This reaction yields an adsorbed species with an associated equilibrium constant :
A_{g}{} + S <=> A_{ad}.
From these basic hypotheses the mathematical formulation of the Langmuir adsorption isotherm can be derived in various independent and complementary ways: by the kinetics, the thermodynamics, and the statistical mechanics approaches respectively (see below for the different demonstrations).
The Langmuir adsorption equation is
where is the fractional occupancy of the adsorption sites, i.e., the ratio of the volume of gas adsorbed onto the solid to the volume of a gas molecules monolayer covering the whole surface of the solid and completely occupied by the adsorbate. A continuous monolayer of adsorbate molecules covering a homogeneous flat solid surface is the conceptual basis for this adsorption model.
Background and experiments
In 1916, Irving Langmuir presented his model for the adsorption of species onto simple surfaces. Langmuir was awarded the Nobel Prize in 1932 for his work concerning surface chemistry. He hypothesized that a given surface has a certain number of equivalent sites to which a species can "stick", either by physisorption or chemisorption. His theory began when he postulated that gaseous molecules do not rebound elastically from a surface, but are held by it in a similar way to groups of molecules in solid bodies.
Langmuir published two papers that confirmed the assumption that adsorbed films do not exceed one molecule in thickness. The first experiment involved observing electron emission from heated filaments in gases. The second, a more direct evidence, examined and measured the films of liquid onto an adsorbent surface layer. He also noted that generally the attractive strength between the surface and the first layer of adsorbed substance is much greater than the strength between the first and second layer. However, there are instances where the subsequent layers may condense given the right combination of temperature and pressure.
Basic assumptions of the model
Inherent within this model, the following assumptions are valid specifically for the simplest case: the adsorption of a single adsorbate onto a series of equivalent sites onto the surface of the solid.
The surface containing the adsorbing sites is a perfectly flat plane with no corrugations (assume the surface is homogeneous). However, chemically heterogeneous surfaces can be considered to be homogeneous if the adsorbate is bound to only one type of functional groups on the surface.
The adsorbing gas adsorbs into an immobile state.
All sites are energetically equivalent, and the energy of adsorption is equal for all sites.
Each site can hold at most one molecule (mono-layer coverage only).
No (or ideal) interactions between adsorbate molecules on adjacent sites. When the interactions are ideal, the energy of side-to-side interactions is equal for all sites regardless of the surface occupancy.
Derivations of the Langmuir adsorption isotherm
The mathematical expression of the Langmuir adsorption isotherm involving only one sorbing species can be demonstrated in different ways: the kinetics approach, the thermodynamics approach, and the statistical mechanics approach respectively. In case of two competing adsorbed species, the competitive adsorption model is required, while when a sorbed species dissociates into two distinct entities, the dissociative adsorption model need to be used.
Kinetic derivation
This section provides a kinetic derivation for a single-adsorbate case. The kinetic derivation applies to gas-phase adsorption. However, it has been mistakenly applied to solutions. The multiple-adsorbate case is covered in the competitive adsorption sub-section.
The model assumes adsorption and desorption as being elementary processes, where the rate of adsorption rad and the rate of desorption rd are given by
where pA is the partial pressure of A over the surface, [S] is the concentration of free sites in number/m2, [Aad] is the surface concentration of A in molecules/m2 (concentration of occupied sites), and kad and kd are constants of forward adsorption reaction and backward desorption reaction in the above reactions.
At equilibrium, the rate of adsorption equals the rate of desorption. Setting rad = rd and rearranging, we obtain
The concentration of sites is given by dividing the total number of sites (S0) covering the whole surface by the area of the adsorbent (a):
We can then calculate the concentration of all sites by summing the concentration of free sites [S] and occupied sites:
Combining this with the equilibrium equation, we get
We define now the fraction of the surface sites covered with A as
This, applied to the previous equation that combined site balance and equilibrium, yields the Langmuir adsorption isotherm:
Thermodynamic derivation
In condensed phases (solutions), adsorption to a solid surface is a competitive process between the solvent (A) and the solute (B) to occupy the binding site. The thermodynamic equilibrium is described as
Solvent (bound) + Solute (free) ↔ Solvent (free) + Solute (bound).
If we designate the solvent by the subscript "1" and the solute by "2", and the bound state by the superscript "s" (surface/bound) and the free state by the "b" (bulk solution / free), then the equilibrium constant can be written as a ratio between the activities of products over reactants:
For dilute solutions the activity of the solvent in bulk solution and the activity coefficients () are also assumed to ideal on the surface. Thus, , and where are mole fractions.
Re-writing the equilibrium constant and solving for yields
Note that the concentration of the solute adsorbate can be used instead of the activity coefficient. However, the equilibrium constant will no longer be dimensionless and will have units of reciprocal concentration instead. The difference between the kinetic and thermodynamic derivations of the Langmuir model is that the thermodynamic uses activities as a starting point while the kinetic derivation uses rates of reaction. The thermodynamic derivation allows for the activity coefficients of adsorbates in their bound and free states to be included. The thermodynamic derivation is usually referred to as the "Langmuir-like equation".
Statistical mechanical derivation
This derivation
based on statistical mechanics was originally provided by Volmer and Mahnert in 1925. The partition function of the finite number of adsorbents adsorbed on a surface, in a canonical ensemble, is given by
where is the partition function of a single adsorbed molecule, is the number of adsorption sites (both occupied and unoccupied), and is the number of adsorbed molecules which should be less than or equal to . The terms in the bracket give the total partition function of the adsorbed molecules by taking a product of the individual partition functions (refer to Partition function of subsystems). The factor accounts for the overcounting arising due to the indistinguishable nature of the adsorbates. The grand canonical partition function is given by
is the chemical potential of an adsorbed molecule. As it has the form of binomial series, the summation is reduced to
where
The grand canonical potential is
based on which the average number of occupied sites is calculated
which gives the coverage
Now, invoking the condition that the system is in equilibrium, that is, the chemical potential of the adsorbed molecules is equal to that of the molecules in gas phase, we have
The chemical potential of an ideal gas is
where is the Helmholtz free energy of an ideal gas with its partition function
is the partition function of a single particle in the volume of (only consider the translational freedom here).
We thus have , where we use Stirling's approximation.
Plugging to the expression of , we have
which gives the coverage
By defining
and using the identity , finally, we have
It is plotted in the figure alongside demonstrating that the surface coverage increases quite rapidly with the partial pressure of the adsorbants, but levels off after P reaches P0.
Competitive adsorption
The previous derivations assumed that there is only one species, A, adsorbing onto the surface. This section considers the case when there are two distinct adsorbates present in the system. Consider two species A and B that compete for the same adsorption sites. The following hypotheses are made here:
All the sites are equivalent.
Each site can hold at most one molecule of A, or one molecule of B, but not both simultaneously.
There are no interactions between adsorbate molecules on adjacent sites.
As derived using kinetic considerations, the equilibrium constants for both A and B are given by
and
The site balance states that the concentration of total sites [S0] is equal to the sum of free sites, sites occupied by A and sites occupied by B:
Inserting the equilibrium equations and rearranging in the same way we did for the single-species adsorption, we get similar expressions for both θA and θB:
Dissociative adsorption
The other case of special importance is when a molecule D2 dissociates into two atoms upon adsorption. Here, the following assumptions would be held to be valid:
D2 completely dissociates to two molecules of D upon adsorption.
The D atoms adsorb onto distinct sites on the surface of the solid and then move around and equilibrate.
All sites are equivalent.
Each site can hold at most one atom of D.
There are no interactions between adsorbate molecules on adjacent sites.
Using similar kinetic considerations, we get
The 1/2 exponent on pD2 arises because one gas phase molecule produces two adsorbed species. Applying the site balance as done above,
Entropic considerations
The formation of Langmuir monolayers by adsorption onto a surface dramatically reduces the entropy of the molecular system.
To find the entropy decrease, we find the entropy of the molecule when in the adsorbed condition.
Using Stirling's approximation, we have
On the other hand, the entropy of a molecule of an ideal gas is
where is the thermal de Broglie wavelength of the gas molecule.
Limitations of the model
The Langmuir adsorption model deviates significantly in many cases, primarily because it fails to account for the surface roughness of the adsorbent. Rough inhomogeneous surfaces have multiple site types available for adsorption, with some parameters varying from site to site, such as the heat of adsorption. Moreover, specific surface area is a scale-dependent quantity, and no single true value exists for this parameter. Thus, the use of alternative probe molecules can often result in different obtained numerical values for surface area, rendering comparison problematic.
The model also ignores adsorbate–adsorbate interactions. Experimentally, there is clear evidence for adsorbate–adsorbate interactions in heat of adsorption data. There are two kinds of adsorbate–adsorbate interactions: direct interaction and indirect interaction. Direct interactions are between adjacent adsorbed molecules, which could make adsorbing near another adsorbate molecule more or less favorable and greatly affects high-coverage behavior. In indirect interactions, the adsorbate changes the surface around the adsorbed site, which in turn affects the adsorption of other adsorbate molecules nearby.
Modifications
The modifications try to account for the points mentioned in above section like surface roughness, inhomogeneity, and adsorbate–adsorbate interactions.
Two-mechanism Langmuir-like equation (TMLLE)
Also known as the two-site Langmuir equation. This equation describes the adsorption of one adsorbate to two or more distinct types of adsorption sites. Each binding site can be described with its own Langmuir expression, as long as the adsorption at each binding site type is independent from the rest.
where
– total amount adsorbed at a given adsorbate concentration,
– maximum capacity of site type 1,
– maximum capacity of site type 2,
– equilibrium (affinity) constant of site type 1,
– equilibrium (affinity) constant of site type 2,
– adsorbate activity in solution at equilibrium
This equation works well for adsorption of some drug molecules to activated carbon in which some adsorbate molecules interact with hydrogen bonding while others interact with a different part of the surface by hydrophobic interactions (hydrophobic effect). The equation was modified to account for the hydrophobic effect (also known as entropy-driven adsorption):
The hydrophobic effect is independent of concentration, since Therefore, the capacity of the adsorbent for hydrophobic interactions can obtained from fitting to experimental data. The entropy-driven adsorption originates from the restriction of translational motion of bulk water molecules by the adsorbate, which is alleviated upon adsorption.
Freundlich adsorption isotherm
The Freundlich isotherm is the most important multi-site adsorption isotherm for rough surfaces.
where αF and CF are fitting parameters. This equation implies that if one makes a log–log plot of adsorption data, the data will fit a straight line. The Freundlich isotherm has two parameters, while Langmuir's equations has only one: as a result, it often fits the data on rough surfaces better than the Langmuir isotherm. However, the Freundlich equation is not unique; consequently, a good fit of the data points does not offer sufficient proof that the surface is heterogeneous. The heterogeneity of the surface can be confirmed with calorimetry. Homogeneous surfaces (or heterogeneous surfaces that exhibit homogeneous adsorption (single-site)) have a constant of adsorption as a function of the occupied-sites fraction. On the other hand, heterogeneous adsorbents (multi-site) have a variable of adsorption depending on the sites occupation. When the adsorbate pressure (or concentration) is low, the fractional occupation is small and as a result, only low-energy sites are occupied, since these are the most stable. As the pressure increases, the higher-energy sites become occupied, resulting in a smaller of adsorption, given that adsorption is an exothermic process.
A related equation is the Toth equation. Rearranging the Langmuir equation, one can obtain
J. Toth modified this equation by adding two parameters αT0 and CT0 to formulate the Toth equation:
Temkin adsorption isotherm
This isotherm takes into account indirect adsorbate–adsorbate interactions on adsorption isotherms. Temkin noted experimentally that heats of adsorption would more often decrease than increase with increasing coverage.
The heat of adsorption ΔHad is defined as
He derived a model assuming that as the surface is loaded up with adsorbate, the heat of adsorption of all the molecules in the layer would decrease linearly with coverage due to adsorbate–adsorbate interactions:
where αT is a fitting parameter. Assuming the Langmuir adsorption isotherm still applied to the adsorbed layer, is expected to vary with coverage as follows:
Langmuir's isotherm can be rearranged to
Substituting the expression of the equilibrium constant and taking the natural logarithm:
BET equation
Brunauer, Emmett and Teller (BET) derived the first isotherm for multilayer adsorption. It assumes a random distribution of sites that are empty or that are covered with by one monolayer, two layers and so on, as illustrated alongside. The main equation of this model is
where
and [A] is the total concentration of molecules on the surface, given by
where
in which [A]0 is the number of bare sites, and [A]i is the number of surface sites covered by i molecules.
Adsorption of a binary liquid on a solid
This section describes the surface coverage when the adsorbate is in liquid phase and is a binary mixture.
For ideal both phases no lateral interactions, homogeneous surface the composition of a surface phase for a binary liquid system in contact with solid surface is given by a classic Everett isotherm equation (being a simple analogue of Langmuir equation), where the components are interchangeable (i.e. "1" may be exchanged to "2") without change of equation form:
where the normal definition of multi-component system is valid as follows:
By simple rearrangement, we get
This equation describes competition of components "1" and "2".
| Physical sciences | Other separations | Chemistry |
1345538 | https://en.wikipedia.org/wiki/Anticline | Anticline | In structural geology, an anticline is a type of fold that is an arch-like shape and has its oldest beds at its core, whereas a syncline is the inverse of an anticline. A typical anticline is convex up in which the hinge or crest is the location where the curvature is greatest, and the limbs are the sides of the fold that dip away from the hinge. Anticlines can be recognized and differentiated from antiforms by a sequence of rock layers that become progressively older toward the center of the fold. Therefore, if age relationships between various rock strata are unknown, the term antiform should be used.
The progressing age of the rock strata towards the core and uplifted center, are the trademark indications for evidence of anticlines on a geologic map. These formations occur because anticlinal ridges typically develop above thrust faults during crustal deformations. The uplifted core of the fold causes compression of strata that preferentially erodes to a deeper stratigraphic level relative to the topographically lower flanks. Motion along the fault including both shortening and extension of tectonic plates, usually also deforms strata near the fault. This can result in an asymmetrical or overturned fold.
Terminology of different folds
Antiform
An antiform can be used to describe any fold that is convex up. It is the relative ages of the rock strata that distinguish anticlines from antiforms.
Elements
The hinge of an anticline refers to the location where the curvature is greatest, also called the crest. The hinge is also the highest point on a stratum along the top of the fold. The culmination also refers to the highest point along any geologic structure. The limbs are the sides of the fold that display less curvature. The inflection point is the area on the limbs where the curvature changes direction.
The axial surface is an imaginary plane connecting the hinge of each layer of rock stratum through the cross section of an anticline. If the axial surface is vertical and the angles on each side of the fold are equivalent, then the anticline is symmetrical. If the axial plane is tilted or offset, then the anticline is asymmetrical. An anticline that is cylindrical has a well-defined axial surface, whereas non-cylindrical anticlines are too complex to have a single axial plane.
Types
An overturned anticline is an asymmetrical anticline with a limb that has been tilted beyond perpendicular, so that the beds in that limb have basically flipped over and may dip in the same direction on both sides of the axial plane. If the angle between the limbs is large (70–120 degrees), then the fold is an "open" fold, but if the angle between the limbs is small (30 degrees or less), then the fold is a "tight" fold. If an anticline plunges (i.e., the anticline crest is inclined to the Earth's surface), it will form Vs on a geologic map view that point in the direction of plunge. A plunging anticline has a hinge that is not parallel to the earth's surface. All anticlines and synclines have some degree of plunge. Periclinal folds are a type of anticlines that have a well-defined, but curved hinge line and are doubly plunging and thus elongate domes.
Folds in which the limbs dip toward the hinge and display a more U-like shape are called synclines. They usually flank the sides of anticlines and display opposite characteristics. A syncline's oldest rock strata are in its outer limbs; the rocks become progressively younger toward its hinge. A monocline is a bend in the strata resulting in a local steepening in only one direction of dip. Monoclines have the shape of a carpet draped over a stairstep.
An anticline that has been more deeply eroded in the center is called a breached or scalped anticline. Breached anticlines can become incised by stream erosion, forming an anticlinal valley.
A structure that plunges in all directions to form a circular or elongate structure is a dome. Domes may be created via diapirism from underlying magmatic intrusions or upwardly mobile, mechanically ductile material such as rock salt (salt dome) and shale (shale diapir) that cause deformations and uplift in the surface rock. The Richat Structure of the Sahara is considered a dome that has been laid bare by erosion.
An anticline which plunges at both ends is termed a doubly plunging anticline, and may be formed from multiple deformations, or superposition of two sets of folds. It may also be related to the geometry of the underlying detachment fault and the varying amount of displacement along the surface of that detachment fault.
An anticlinorium is a large anticline in which a series of minor anticlinal folds are superimposed. Examples include the Late Jurassic to Early Cretaceous Purcell Anticlinorium in British Columbia and the Blue Ridge anticlinorium of northern Virginia and Maryland in the Appalachians, or the Nittany Valley in central Pennsylvania.
Formation processes
Anticlines are usually developed above thrust faults, so any small compression and motion within the inner crust can have large effects on the upper rock stratum. Stresses developed during mountain building or during other tectonic processes can similarly warp or bend bedding and foliation (or other planar features). The more the underlying fault is tectonically uplifted, the more the strata will be deformed and must adapt to new shapes. The shape formed will also be very dependent on the properties and cohesion of the different types of rock within each layer.
During the formation of flexural-slip folds, the different rock layers form parallel-slip folds to accommodate for buckling. A good way to visualize how the multiple layers are manipulated, is to bend a deck of cards and to imagine each card as a layer of rock stratum. The amount of slip on each side of the anticline increases from the hinge to the inflection point.
Passive-flow folds form when the rock is so soft that it behaves like weak plastic and slowly flows. In this process different parts of the rock body move at different rates causing shear stress to gradually shift from layer to layer. There is no mechanical contrast between layers in this type of fold. Passive-flow folds are extremely dependent on the rock composition of the stratum and can typically occur in areas with high temperatures.
Economic significance
Anticlines, structural domes, fault zones and stratigraphic traps are very favorable locations for oil and natural gas drilling. About 80 percent of the world's petroleum has been found in anticlinal traps. The low density of petroleum causes oil to buoyantly migrate out of its source rock and upward toward the surface until it is trapped and stored in reservoir rock such as sandstone or porous limestone. The oil becomes trapped along with water and natural gas by a caprock that is made up of impermeable barrier such as an impermeable stratum or fault zone. Examples of low-permeability seals that contain the hydrocarbons, oil and gas, in the ground include shale, limestone, sandstone, and rock salt. The actual type of stratum does not matter as long as it has low permeability.
Water, minerals and specific rock strata such as limestone found inside anticlines are also extracted and commercialized. Lastly, ancient fossils are often found in anticlines and are used for paleontological research or harvested into products to be sold.
Notable examples
Asia
Ghawar Anticline, Saudi Arabia, the structural trap for the largest conventional oil field in the world.
Australia
Hill End Anticline, New South Wales, which is associated with deposits of gold.
Castlemaine Anticlinal Fold, Victoria, which is celebrated with a plaque that says, "This fine exhibit was disclosed when Lyttleton Street East was constructed in 1874. Saddle reefs occur in similar folds of the sandstones and slates on lower geological horizons."
Europe
The Weald–Artois Anticline is a major anticline which outcrops in southeast England and northern France. It was formed from the late Oligocene to middle Miocene, during the Alpine orogeny.
North America
Anticlines can have a major effect on the local geomorphology and economy of the regions in which they occur. One example of this is the El Dorado anticline in Kansas. The anticline was first tapped into for its petroleum in 1918. Soon after the site became a very prosperous area for entrepreneurs following World War I and the rapid popularization of motor vehicles. By 1995 the El Dorado oil fields had produced 300 million barrels of oil. The central Kansas uplift is an antiform composed of several small anticlines that have collectively produced more than 2.5 million barrels of oil.
Another notable anticline is the Tierra Amarilla anticline in San Ysidro, New Mexico. This is a popular hiking and biking site because of the great biodiversity, geologic beauty and paleontological resources. This plunging anticline is made up of Petrified Forest mudstones and sandstone and its caprock is made of Pleistocene and Holocene travertine. The anticline contains springs that deposit carbon dioxide travertine that help to contribute to the rich diversity of microorganisms. This area also contains remains of fossils and ancient plants from the Jurassic period that are sometimes exposed through geological erosion.
The Ventura Anticline is a geologic structure that is part of the Ventura oil fields, the seventh largest oil field in California that was discovered in the 1860s. The anticline runs east to west for 16 miles, dipping steeply 30–60 degrees at both ends. Ventura County has a high rate of compression and seismic activity due to the converging San Andreas Fault. As a result, the Ventura anticline rises at a rate of 5 mm/year with the adjacent Ventura Basin converging at a rate of about 7–10 mm/year. The anticline is composed of a series of sandstone rock beds and an impermeable rock cap under which vast reserves of oil and gas are trapped. Eight different oil bearing zones along the anticline vary greatly from 3,500 to 12,000 feet. The oil and gas formed these pools as they migrated upward during the Pliocene Era and became contained beneath the caprock. This oil field is still active and has a cumulative production of one billion barrels of oil making it one of the most vital historical and economic features of Ventura County.
Gallery
| Physical sciences | Structural geology | Earth science |
1345771 | https://en.wikipedia.org/wiki/Position%20%28geometry%29 | Position (geometry) | In geometry, a position or position vector, also known as location vector or radius vector, is a Euclidean vector that represents a point P in space. Its length represents the distance in relation to an arbitrary reference origin O, and its direction represents the angular orientation with respect to given reference axes. Usually denoted x, r, or s, it corresponds to the straight line segment from O to P.
In other words, it is the displacement or translation that maps the origin to P:
The term position vector is used mostly in the fields of differential geometry, mechanics and occasionally vector calculus.
Frequently this is used in two-dimensional or three-dimensional space, but can be easily generalized to Euclidean spaces and affine spaces of any dimension.
Relative position
The relative position of a point Q with respect to point P is the Euclidean vector resulting from the subtraction of the two absolute position vectors (each with respect to the origin):
where .
The relative direction between two points is their relative position normalized as a unit vector
Definition and representation
Three dimensions
In three dimensions, any set of three-dimensional coordinates and their corresponding basis vectors can be used to define the location of a point in space—whichever is the simplest for the task at hand may be used.
Commonly, one uses the familiar Cartesian coordinate system, or sometimes spherical polar coordinates, or cylindrical coordinates:
where t is a parameter, owing to their rectangular or circular symmetry. These different coordinates and corresponding basis vectors represent the same position vector. More general curvilinear coordinates could be used instead and are in contexts like continuum mechanics and general relativity (in the latter case one needs an additional time coordinate).
n dimensions
Linear algebra allows for the abstraction of an n-dimensional position vector. A position vector can be expressed as a linear combination of basis vectors:
The set of all position vectors forms position space (a vector space whose elements are the position vectors), since positions can be added (vector addition) and scaled in length (scalar multiplication) to obtain another position vector in the space. The notion of "space" is intuitive, since each xi (i = 1, 2, …, n) can have any value, the collection of values defines a point in space.
The dimension of the position space is n (also denoted dim(R) = n). The coordinates of the vector r with respect to the basis vectors ei are xi. The vector of coordinates forms the coordinate vector or n-tuple (x1, x2, …, xn).
Each coordinate xi may be parameterized a number of parameters t. One parameter xi(t) would describe a curved 1D path, two parameters xi(t1, t2) describes a curved 2D surface, three xi(t1, t2, t3) describes a curved 3D volume of space, and so on.
The linear span of a basis set B = {e1, e2, …, en} equals the position space R, denoted span(B) = R.
Applications
Differential geometry
Position vector fields are used to describe continuous and differentiable space curves, in which case the independent parameter needs not be time, but can be (e.g.) arc length of the curve.
Mechanics
In any equation of motion, the position vector r(t) is usually the most sought-after quantity because this function defines the motion of a particle (i.e. a point mass) – its location relative to a given coordinate system at some time t.
To define motion in terms of position, each coordinate may be parametrized by time; since each successive value of time corresponds to a sequence of successive spatial locations given by the coordinates, the continuum limit of many successive locations is a path the particle traces.
In the case of one dimension, the position has only one component, so it effectively degenerates to a scalar coordinate. It could be, say, a vector in the x direction, or the radial r direction. Equivalent notations include
Derivatives
For a position vector r that is a function of time t, the time derivatives can be computed with respect to t. These derivatives have common utility in the study of kinematics, control theory, engineering and other sciences.
Velocity
where dr is an infinitesimally small displacement (vector).
Acceleration
Jerk
These names for the first, second and third derivative of position are commonly used in basic kinematics. By extension, the higher-order derivatives can be computed in a similar fashion. Study of these higher-order derivatives can improve approximations of the original displacement function. Such higher-order terms are required in order to accurately represent the displacement function as a sum of an infinite sequence, enabling several analytical techniques in engineering and physics.
| Physical sciences | Classical mechanics | Physics |
1345963 | https://en.wikipedia.org/wiki/Sumatran%20tiger | Sumatran tiger | The Sumatran tiger is a population of Panthera tigris sondaica on the Indonesian island of Sumatra. It is the only surviving tiger population in the Sunda Islands, where the Bali and Javan tigers are extinct.
Sequences from complete mitochondrial genes of 34 tigers support the hypothesis that Sumatran tigers are diagnostically distinct from mainland subspecies. In 2017, the Cat Classification Task Force of the Cat Specialist Group revised felid taxonomy and recognizes the living and extinct tiger populations in Indonesia as P. t. sondaica.
Taxonomy
Felis tigris sondaicus was the scientific name proposed by Coenraad Jacob Temminck in 1844 for a tiger specimen from Java.
Panthera tigris sumatrae was proposed by Reginald Innes Pocock in 1929, who described a skin and a skull of a tiger zoological specimen from Sumatra.
The skull and pelage pattern of tiger specimens from Java and Sumatra do not differ significantly.
P. t. sondaica is therefore considered the valid name for the living and extinct tiger populations in Indonesia.
Evolution
Analysis of DNA is consistent with the hypothesis that Sumatran tigers became isolated from other tiger populations after a rise in sea level that occurred at the Pleistocene to Holocene transition about 12,000–6,000 years ago. In agreement with this evolutionary history, the Sumatran tiger is genetically isolated from all living mainland tigers, which form a distinct group closely related to each other. The isolation of the Sumatran tiger from mainland tiger populations is supported by multiple unique characters, including two diagnostic mitochondrial DNA nucleotide sites, ten mitochondrial DNA haplotypes and 11 out of 108 unique microsatellite alleles. The relatively high genetic variability and the phylogenetic distinctiveness of the Sumatran tiger indicates that the gene flow between island and mainland populations was highly restricted.
Characteristics
The Sumatran tiger was described based on two zoological specimens that differed in skull size and striping pattern from Bengal and Javan tiger specimens. It is darker in fur colour and has broader stripes than the Javan tiger. Stripes tend to dissolve into spots near their ends, and on the back, flanks and hind legs are lines of small, dark spots between the regular stripes. The frequency of stripes is higher than in other subspecies.
Males have a prominent ruff, which is especially marked in the Sumatran tiger.
The Sumatran tiger is one of the smallest tigers. Males measure between the pegs in head-to-body length, with the greatest skull length of and weigh . Females weigh and measure in length between the pegs with a greatest length of skull of .
Distribution and habitat
The Sumatran tiger persists in small and fragmented populations across Sumatra, from sea level in the coastal lowland forest of Bukit Barisan Selatan National Park on the southeastern tip of Lampung Province to in mountain forests of Gunung Leuser National Park in Aceh Province. It is present in 27 habitat patches larger than , which cover . About a third of these patches are inside protected areas.
Sumatran tigers prefer lowland and hill forests, where up to three tigers live in an area of ; they use non-forest habitats and human-dominated landscapes at the fringes of protected areas to a lesser degree.
In 1978, the Sumatran tiger population was estimated at 1,000 individuals, based on responses to a questionnaire survey. In 1985, a total of 26 protected areas across Sumatra containing about 800 tigers were identified. In 1992, an estimated 400–500 tigers lived in five Sumatran national parks and two protected areas. At that time, the largest population unit comprised 110–180 individuals in Gunung Leuser National Park.
As of 2011, the tiger population in Kerinci Seblat National Park in central Sumatra comprised 165–190 individuals, which is more than anywhere else on the island. The park has the highest tiger occupancy rate of Sumatra's protected areas, with 83% of the park showing signs of tigers.
Sumatra's total tiger population was estimated at 618 ± 290 individuals in 2017.
Ecology and behaviour
Sumatran tigers strongly prefer uncultivated forests and make little use of plantations of acacia and oil palm even if these are available. Within natural forest areas, they tend to use areas with higher elevation, lower annual rainfall, farther from the forest edge, and closer to forest centres. They prefer forest with dense understory cover and steep slope, and they strongly avoid forest areas with high human influence in the forms of encroachment and settlement. In acacia plantations, they tend to use areas closer to water and prefer areas with older plants, more leaf litter, and thicker subcanopy cover. Tiger records in oil palm plantations and rubber plantations are scarce. The availability of adequate vegetation cover at the ground level serves as an environmental condition fundamentally needed by tigers regardless of the location. Without adequate understory cover, tigers are even more vulnerable to persecution by humans. Human disturbance-related variables negatively affect tiger occupancy and habitat use. Variables with strong impacts include settlement and encroachment within forest areas, logging, and the intensity of maintenance in acacia plantations.
Camera trapping surveys conducted in southern Riau revealed an extremely low abundance of potential prey and a low tiger density in peat swamp forest areas. Repeated sampling in the newly established Tesso Nilo National Park documented a trend of increasing tiger density from 0.90 individuals per in 2005 to 1.70 individuals per in 2008.
In the Bukit Barisan Selatan National Park, nine prey species larger than of body weight were identified including great argus, pig-tailed macaque, Malayan porcupine, Malayan tapir, banded pig, greater and lesser mouse-deer, Indian muntjac, and Sambar deer.
Threats
Major threats include habitat loss due to expansion of palm oil plantations and planting of acacia plantations, prey-base depletion, and illegal trade primarily for the domestic market.
Conflicts with humans are another major threat to the Sumatran tiger. Poachers target tigers with wire snares, and they are also inadvertently caught in traps set by deer hunters and farmers attempting to control crop raids from wild boar.
Tigers need large contiguous forest blocks to thrive. Between 1985 and 1999, forest loss within Bukit Barisan Selatan National Park averaged 2% per year. A total of of forest disappeared inside the park, and were lost in a 10-km buffer, eliminating forest outside the park. Lowland forest disappeared faster than montane forest, and forests on gentle slopes disappeared faster than forests on steep slopes. Most forest conversions resulted from agricultural development, leading to predictions that by 2010, 70% of the park will be in agriculture. Camera-trap data indicated avoidance of forest boundaries by tigers. Classification of forest into core and peripheral forest based on mammal distribution suggests that by 2010, core forest area for tigers will be fragmented and reduced to 20% of the remaining forest.
Sumatra's largest tiger population in Kerinci Seblat National Park is threatened by a high rate of deforestation in its outer regions. Drivers are an unsustainable demand for natural resources created by a human population with the highest rate of growth in Indonesia, and a government initiative to increase tree-crop plantations and high-intensity commercial logging, which ultimately leads to forest fires. The majority of the tigers found in the park were relocated to its center where conservation efforts are focused, but issues in the lowland hill forests of the outskirts remain. While being a highly suitable tiger habitat, these areas are also heavily targeted by logging efforts, which substantially contributes to declines in local tiger numbers.
The expansion of plantations is increasing greenhouse gas emissions, playing a part in anthropogenic climate change, thus further adding to environmental pressures on endangered species. Climate-based movement of tigers northwards may lead to increased conflict with people. From 1987 to 1997, Sumatran tigers reportedly killed 146 people and at least 870 livestock. In West Sumatra, Riau, and Aceh, a total of 128 incidents were reported; 265 tigers were killed and 97 captured in response, and 35 more tigers were killed from 1998 to 2002. From 2007 to 2010, the tigers caused the death of 9 humans and 25 further tigers were killed.
In 1997, an estimated 53 tigers were killed by poachers and their parts sold throughout most of northern Sumatra. Numbers for all of Sumatra are likely to be higher. Farmers killed many of the tigers to prevent livestock losses. They sold them to gold and souvenir shops, and pharmacies. In 2006, wildlife markets were surveyed in 28 cities and nine seaports in seven Sumatran provinces; 33 of 326 retail outlets offered tiger parts like skins, canines, bones, and whiskers. Tiger bones fetched the highest average price of US$116 per kg, followed by canines. There is evidence that tiger parts are smuggled out of Indonesia. In July 2005, over of tiger bones and 24 skulls were confiscated in Taiwan in a shipment from Jakarta.
In 2013–2014, Kerinci Seblat National Park experienced an upsurge in poaching, with the highest annual number of snare traps being removed for a patrol effort similar to previous years. Evidence is scarce and misunderstood on whether the strategies implemented to diminish poaching are succeeding despite the investment of millions of dollars annually into conservation strategies.
In provincially-managed forests in Aceh province, Sumatran tigers are threatened by poaching due to insufficient or nonexistent ranger patrols.
Conservation
Panthera tigris is listed on CITES Appendix I. Hunting is prohibited in Indonesia.
In 1994, the Indonesian Sumatran Tiger Conservation Strategy addressed the potential crisis that tigers faced in Sumatra. The Sumatran Tiger Project (STP) was initiated in June 1995 in and around the Way Kambas National Park to ensure the long-term viability of wild Sumatran tigers and to accumulate data on tiger life-history characteristics vital for the management of wild populations. By August 1999, the teams of the STP had evaluated 52 sites of potential tiger habitat in Lampung Province, of which only 15 were intact enough to contain tigers. In the framework of the STP, a community-based conservation programme was initiated to document the tiger-human dimension in the park to enable conservation authorities to resolve tiger-human conflicts based on a comprehensive database rather than anecdotes and opinions.
In 2007, the Indonesian Forestry Ministry and Safari Park established cooperation with the Australia Zoo for the conservation of Sumatran tigers and other endangered species. The program includes conserving Sumatran tigers and other endangered species in the wild, efforts to reduce conflicts between tigers and humans, and rehabilitating Sumatran tigers and reintroducing them to their natural habitat.
Indonesia's struggle with conservation has caused an upsurge in political momentum to protect and conserve wildlife and biodiversity. In 2009, Indonesia's president committed to substantially reduce deforestation, and policies across the nation requiring spatial plans that would be environmentally sustainable at national, provincial, and district levels.
Between 2005 and 2015, about US$210 million have been invested into tiger law-enforcement activities that support forest ranger patrols, as well as the implementations of front-line law-enforcement activities by the Global Tiger Recovery Plan, which aims to double the number of wild tigers by 2020.
In November 2016, Batu Nanggar Sanctuary was opened in North Padang Lawas Regency, North Sumatra
for conservation of Sumatran wildlife.
An interview survey among 600 consumers revealed that most were willing to pay consistently more for a "tiger-friendly" produced good if this product would be conducive to Sumatran tiger conservation.
In captivity
As of 2013, about 375 captive Sumatran tigers were listed in the global studbook and management plan, with 50 of them housed in 14 zoos in Australia and New Zealand. All of them were offspring of 15 founders. Fourteen cubs showed congenital vestibular system dysfunctions such as ataxia, strabismus, nystagmus, head tilting and falling that resolved when they were two years old. The cause for this disorder is most likely their close genetic relation and inbreeding.
| Biology and health sciences | Felines | Animals |
1346982 | https://en.wikipedia.org/wiki/Gharial | Gharial | The gharial (Gavialis gangeticus), also known as gavial or fish-eating crocodile, is a crocodilian in the family Gavialidae and among the longest of all living crocodilians. Mature females are long, and males . Adult males have a distinct boss at the end of the snout, which resembles an earthenware pot known as a ghara, hence the name "gharial". The gharial is well adapted to catching fish because of its long, narrow snout and 110 sharp, interlocking teeth.
The gharial probably evolved in the northern Indian subcontinent. Fossil gharial remains were excavated in Pliocene deposits in the Sivalik Hills and the Narmada River valley. It currently inhabits rivers in the plains of the northern part of the Indian subcontinent. It is the most thoroughly aquatic crocodilian, and leaves the water only for basking and building nests on moist sandbanks. Adults mate at the end of the cold season. Females congregate in spring to dig nests, in which they lay 20–95 eggs. They guard the nests and the young, which hatch before the onset of the monsoon. The hatchlings stay and forage in shallow water during their first year, but move to sites with deeper water as they grow.
The wild gharial population has declined drastically since the 1930s and is limited to only 2% of its historical range today. Conservation programmes initiated in India and Nepal focused on reintroducing captive-bred gharials since the early 1980s. Loss of habitat because of sand mining and conversion to agriculture, depletion of fish resources and detrimental fishing methods continue to threaten the population. It has been listed as critically endangered on the IUCN Red List since 2007.
The oldest known depictions of the gharial are about 4,000 years old and were found in the Indus Valley. Hindus regard it as the vehicle of the river deity Gaṅgā. Local people living near rivers attributed mystical and healing powers to the gharial, and used some of its body parts as ingredients of indigenous medicine.
Etymology
The name 'gharial' is derived from the Hindustani word 'ghara' for an earthen pot, in reference to the nasal protuberance on the adult male's snout. It is also called 'gavial'. The name 'fish-eating crocodile' is a translation of its Bengali name 'mecho kumhir', with 'mecho' being derived from 'māch' meaning fish and 'kumhir' meaning crocodile. The name 'Indian gharial' has occasionally been used for gharial populations in India.
Taxonomy
Lacerta gangetica was the scientific name proposed by Johann Friedrich Gmelin in 1789. Gmelin followed Carl Linnaeus who proposed Lacerta in 1758 to include other crocodiles and various lizards known at the time.
The gharial was placed in the genus Crocodilus by subsequent naturalists:
Crocodilus gavial by Pierre Joseph Bonnaterre in 1789.
Crocodilus longirostris by Johann Gottlob Theaenus Schneider in 1801.
Crocodilus arctirostris by François Marie Daudin in 1802.
Longirostres was a subgroup proposed by Georges Cuvier in 1807 for crocodiles with a long snout. He placed Crocodilus gangeticus with the type locality "Ganges" and Crocodilus tenuirostris without locality into this group.
The generic name Gavialis was proposed by Nicolaus Michael Oppel in 1811 for crocodiles with a cylindrical-shaped back. He placed this genus in the family Crocodilini.
Rhamphostoma was proposed by Johann Georg Wagler in 1830 who considered this genus to contain two species, Crocodilus gangeticus and C. tenuirostris.
The family name Gavialidae was proposed by Arthur Adams in 1854 with Gavialis as the only genus in this family.
Gavialis gangetica was the scientific name used by Albert Günther in 1864 who considered L. gangetica, C. longirostris and C. tenuirostris as synonyms and Gavialis a monotypic taxon. John Edward Gray reviewed zoological specimens in the collection of the Natural History Museum, London. He also considered the gharial monotypic in 1869. He placed it in the family Gavialidae together with the false gharial (Tomistoma schlegelii) because both have long, slender jaws and similar dentition.
Gharialis hysudricus proposed by Richard Lydekker in 1886 was based on a fossil skull from the Sivalik Hills that was larger than gharial fossil skulls known at the time. This name is considered to be a junior synonym of Gavialis gangeticus.
Evolution
The evolution of the gharial and its relationship with and divergence from other crocodilians have been a subject of controversy. Some authors assumed that the gharial evolved earlier than the other crocodilians because of its distinct skull shape and dentition, indicating a more advanced level of specialization. Others suggested that it evolved much later than other crocodilians because of its low levels of blood protein divergence. As it shares this trait with the false gharial, it was suggested that they form a sister group. In contrast, it was suggested that the gharial and all the other crocodilians form a sister group as the structure of its tail muscles is unique. Sequencing of a ribosomal segment of mitochondrial DNAs of gharial and false gharial revealed that they share 22 unique nucleotides, a similarity of 94%, supporting the view that they are sister taxa. Analyses of nuclear gene sequences of both species also support the view that they are sister taxa. Molecular genetics and tip dating studies indicates a genetic divergence between the gharial and false gharial in the Eocene about .
The genus Gavialis probably originated in the region of India and Pakistan in the Early Miocene. Fossil gharial remains excavated in the Sivalik Hills of Haryana and Himachal Pradesh are dated to between the Pliocene and the Early Pleistocene. Fossil gharial remains were also found at two sites in the Ayeyarwady River valley in central Myanmar, which are dated to the Late Pleistocene.
During the Quaternary, Gavialis dispersed as far as Java via the Siva–Malayan route, which did not require saltwater crossings. Fossil remains of Gavialis bengawanicus found on Java were dated to the Early Pleistocene.
G. bengawanicus fossils found in Thailand's Nakhon Ratchasima Province support the hypothesis of gharial dispersal through riverine systems. It represents the only valid extinct Gavialis species.
Phylogeny
The below cladogram of the major extant crocodile groups is based on the latest molecular studies and shows the gharial's close relationship to the false gharial, and how the gavialids and crocodiles are more closely related than the alligatoroids:
Here is a more detailed cladogram that shows the gharial's proposed placement within Gavialidae, including extinct members:
Characteristics
The gharial is olive-coloured, with adults being darker than young, which have dark brown cross bands and speckles. Its back turns almost black at 20 years of age, but its belly is yellowish-white. It has four transverse rows of two scales on the neck, which continue along the back. Scutes on the head, neck and back form a single continuous plate composed of 21 to 22 transverse series, and four longitudinal series. Scutes on the back are bony, but softer and feebly keeled on the sides. The outer edges of the forearms, legs, and feet have crests jutting out; fingers and toes are partly webbed.
Its snout is very long and narrow, widened at the end, and with 27 to 29 upper teeth and 25 or 26 lower teeth on each side. The front teeth are the largest. The first, second, and third lower jaw teeth fit into spaces in the upper jaw. The extremely long mandibular symphysis extends to the 23rd or 24th tooth. The snout of adult gharials is 3.5 times longer than the width of the skull's base. Because of this long snout the gharial is especially adapted to catching and eating fish. The nasal bones are rather short and widely spaced from the premaxillae. The jugal bone is raised. It becomes proportionally thicker with age.
Two individuals in the weight range of had an average measured bite force of .
Male gharials develop a hollow bulbous nasal protuberance at the tip of the snout upon reaching sexual maturity. This protuberance resembles an earthen pot known locally as "ghara". The male's ghara starts growing over the nostrils at an age of 11.5 years and measures about at an age of 15.5 years. It enables the males to emit a hissing sound that can be heard away.
The gharial is the only living crocodilian with such visible sexual dimorphism.
Mature male gharials have larger skulls than females, exceeding a basal length of and a width of .
Female gharials reach sexual maturity at a body length of and grow up to . Males mature at a body length of at least and grow up to a length of . Adult males weigh about on average, but can reach a weight of up to . The gharial is among the largest living crocodilians, with the heaviest recorded male weighing .
A long gharial was claimed to have been killed in the Ghaghara River in Faizabad in August 1920, though no reliable measurements were taken. Male gharials with an alleged length of were sighted around the turn of the 20th century in Indian rivers. Overall, the gharial is less massive when compared to other crocodilians of similar length; a long gharial weighed around , while a long Nile crocodile weighed .
Distribution and habitat
The gharial once thrived in all the major river systems of the northern Indian subcontinent, from the Indus River in Pakistan, the Ganges in India, the Brahmaputra River in northeastern India and Bangladesh to the Irrawaddy River in Myanmar.
In the early 20th century, it was considered common in the Indus River and its Punjabi tributaries. By the early 1980s, it was almost extinct in the Indus. During surveys in 2008 and 2009, no gharial was sighted in the river.
It was also present in India's Godavari River but was hunted to extinction between the late 1940s and the 1960s.
It was considered extinct in the Koshi River since 1970.
In the 1940s, it was numerous in the Barak River in Assam, which held big fish at the time including golden mahseer (Tor putitora).
A few individuals were also sighted in tributaries of the Barak River in Assam, Mizoram and Manipur up to 1988, but surveys were not carried out.
In 1927, a gharial was shot in the Shweli River in Myanmar, a tributary of the Ayeyawady River. This is the only authenticated record in the country attesting the survival of gharials into the 20th century. Whether gharials still live in the Shweli River today is possible but remained unclear in 2012.
By 1976, its global range had decreased to only 2% of its historical range, and fewer than 200 gharials were estimated to survive. It is locally extinct in Pakistan, Bhutan and Myanmar.
Since the early 1980s, the population has been reinforced with captive-bred gharials that were released into wild habitats in India and Nepal. In 2017, the global population was estimated to comprise at maximum 900 individuals, including about 600 mature adults in six major subpopulations along of river courses and another 50 mature adults in eight minor subpopulations along of river courses.
In Nepal, small populations are present and slowly recovering in tributaries of the Ganges, such as the Karnali–Babai River system in Bardia National Park and the Narayani–Rapti river system in Chitwan National Park.
In spring 2017, the Babai River was surveyed using an unmanned aerial vehicle, which detected 33 gharials on a stretch of .
In India, gharial populations are present in the:
Ramganga River in Corbett National Park, where five gharials were recorded in 1974. Captive-bred gharials were released since the late 1970s. The population is breeding since 2008, and increased to about 42 adults by 2013. Most of them congregate along an long stretch of the Kalagarh Reservoir's shoreline. Surveys in 2015 revealed a population of 90 gharials including 59 breeding adults.
Ganges, where 494 gharials were released between 2009 and 2012 in Hastinapur Wildlife Sanctuary.
Girwa River in Katarniaghat Wildlife Sanctuary where the small breeding population was reinforced with captive reared gharials since 1979. A total of 909 gharials were released until 2006, but only 16 nesting females were recorded in the same year. In December 2008, 105 individuals were counted including 35 adults. In spring 2009, 27 nests were detected in seven sites. The number of nest sites decreased from seven in 2017 to two in 2019, possibly due to the upgrowth of woody vegetation and reduced river flow near sandbanks.
Gandaki River downstream the Triveni barrage west of Valmiki Tiger Reserve and adjacent to Sohagi Barwa Sanctuary. The population increased from 15 gharials in 2010 to 54 individuals recorded in March 2015 on a stretch of . 35 of these gharials were wild-born.
Chambal River in National Chambal Sanctuary where 107 gharials were recorded in 1974. Captive-bred gharials were released since 1979, and the population increased to 1,095 gharials in 1992. Between December 2007 and March 2008, 111 gharials were found dead. A total of 948 gharials were counted during surveys in 2013 along the protected river stretch of . In 2017, this population was estimated at 617–761 mature individuals and more than 1250 individuals by two different survey teams; 411 nests were found.
Parbati River, a tributary of the Chambal River, where gharials started using a few sand banks since about 2015; 29 gharials were observed in 2016 and 251 hatchlings counted at two nesting sites in 2017.
Yamuna River where eight young gharials were detected in autumn 2012 near the confluence of the Ken and Yamuna Rivers. They were probably offspring of the breeding population in the Chambal River and had drifted downriver during monsoon floods.
Son River where 164 captive-reared gharials were released between 1981 and 2011.
Koshi River in Bihar where two gharials were sighted basking in late January 2019 during a survey targeting South Asian River Dolphins (Platanista gangetica) on a stretch of about . This is the first record of wild gharials in the river since the 1970s.
Mahanadi River in Odisha's Satkosia Gorge Sanctuary where about 700 gharials were released between 1977 and the early 1990s. During a 1.5 year long survey in 2005–2006, only one male and one female gharial were detected moving together and sharing sand banks in the river.
Between 1979 and 1993, less than 20 individuals were sighted in the upper reaches of the Brahmaputra River between Kaziranga National Park and Dibru-Saikhowa National Park. This population had declined due to commercial fishing, poaching, encroachment by local people in gharial breeding grounds and siltation of river beds following deforestation. In 1998, it was not considered to be viable. About 30 gharials were observed in small lakes and tributaries of the Brahmaputra River in Assam between 2004 and 2007.
In Bangladesh, gharials were recorded in Padma, Jamuna, Mahananda and Brahmaputra rivers between 2000 and 2015.
Behaviour and ecology
The gharial is the most thoroughly aquatic crocodilian. It leaves the water only for basking on riverbanks. Being cold-blooded, it seeks to cool down during hot times and to warm up when ambient temperature is cool. Gharials bask daily in the cold season, foremost in the mornings, and prefer sandy and moist beaches. They change their basking pattern with increasing daily temperatures; they start basking earlier in the mornings, move back into the river when it is hot, and return to the beach later in the afternoon. Groups comprising an adult male, several females and subadults have been observed to bask together. Adult males dominate groups and tolerate immature males.
Large groups of young, subadult and adult gharials form in December and January to bask. Adult males and females associate by mid February.
The gharial shares riverine habitat with the mugger crocodile (Crocodylus palustris) in parts of its range. They use the same nesting grounds, but differ in the selection of basking sites.
The gharial basks close to water on shallow, sandy beaches and lays eggs only in sandy soil near water. The mugger crocodile also basks on sandy beaches, but unlike the gharial climbs steep embankments and rocks, and moves farther away from beaches for both basking and nest building. It also preys on fish, but has a broader prey base than the gharial including snakes, turtles, birds, mammals and dead animals.
Feeding ecology
The gharial is well adapted to hunting fish underwater because of its sharp interlocking teeth and long narrow snout, which meets little resistance in the water. It does not chew its prey, but swallows it whole. Juvenile gharials were observed to jerk their heads back to manoeuvre fish into their gullets, sliding them in head first. Young gharials feed on insects, tadpoles, small fish and frogs. Adults also feed on small crustaceans. Remains of Indian softshell turtle (Nilssonia gangetica) were also found in gharial stomachs. Gharials tear apart large fish and pick up and swallow stones as gastroliths, probably to aid digestion or regulate buoyancy. Some gharial stomachs also contained jewellery.
Stones weighing about were found in a gharial's stomach that was shot in the Sharda River in 1910.
Reproduction
Females mature at a body length of around . Captive females breed at a body length of . Male gharials mature at 15–18 years of age, when they reach a body length of around and once the ghara is developed. The ghara is apparently used to indicate sexual maturity, as a sound resonator when bubbling underwater or for other sexual behaviours.
Courting and mating starts by mid-February at the end of the cold season. In the dry season, reproductive females observed in the Chambal River routinely move and join female breeding groups to dig nests together. They select sites in riverside sand or silt banks located between away from the water and above a water level of . These nests are deep with a diameter of about . Between end of March and early April, they lay 20–95 eggs. A record clutch with 97 eggs was found in Katarniaghat Wildlife Sanctuary.
The eggs are the largest of all crocodilians and weigh an average of . Each egg is long and wide. After 71 to 93 days of incubation, young gharials hatch in July just before the onset of the monsoon. Their sex is most likely determined by temperature, like in most reptiles. Females dig up the hatchlings in response to hatching chirps, but do not assist them to reach the water. They stay at nesting sites until monsoon floods arrive and return after monsoon.
Captive male gharials observed in the 1980s did not participate in guarding nests. A captive male gharial was observed to show an interest in hatchlings and was allowed by the female to carry hatchlings on his back.
In the Chambal River, females were observed to stay close to nest sites and guard young gharials until the shoreline was flooded. VHF radio tracking of a junior male gharial revealed that he was the dominant male guarding nests at a communal nesting site for two years.
Development
Hatchlings range from in body length with a weight of . In two years, they grow to a length of and of in three years.
Gharials hatched and raised in Nepal's Gharial Conservation and Breeding Center measured and weighed at the age of 45 months in April 2013. They consumed up to of fish per individual and month. By the age of 75 months, they had gained in weight and grown reaching body lengths of .
Young gharials in their first year of age hide and forage in shallow water preferably in sites that are surrounded by debris of fallen trees. A study along a stretch of the Chambal River revealed that juvenile gharials up to a body length of prefer basking sites where the mid river water is deep. As their body size increases, they move to sites with deeper water. Subadult and adult gharials above a body length of prefer sites where the water is deeper than .
Young gharials move forward by pushing the diagonally opposite legs synchronously. At a young age, they can also gallop but do so only in emergency situations. When they reach a length of about and a weight of about at the age of 8–9 months, they change to an adult pattern of locomotion of pushing forward with hind and front legs simultaneously. Adults do not have the ability to walk on land in the semi-upright stance as other crocodilians. When basking on the beach, they often turn round so as to face the water.
Threats
The gharial population is estimated to have declined from 5,000–10,000 individuals in 1946 to fewer than 250 individuals in 2006, a decline of 96–98% within three generations. Gharials were killed by fishermen, hunted for skins, trophies and indigenous medicine, and their eggs collected for consumption. The remaining individuals form several fragmented subpopulations. Hunting is no longer considered a significant threat. However, the wild population declined from an estimated 436 adult gharials in 1997 to fewer than 250 mature individuals in 2006. One reason for this decline is the increased use of gill nets for fishing in gharial habitat. The other major reason is the loss of riverine habitat as dams, barrages, irrigation canals and artificial embankments were built; siltation and sand-mining changed river courses; and land near rivers is used for agriculture and grazing by livestock.
When 111 dead gharials were found in the Chambal River between December 2007 and March 2008, it was initially suspected that they had died either because of toxicants or the illegal use of fish nets, in which they became entrapped in and subsequently drowned. Later post mortem pathological testing of tissue samples revealed high levels of heavy metals such as lead and cadmium, which together with stomach ulcers and protozoan parasites reported in most necropsies were thought to have caused their deaths. Water pumps used for pumping water out of the Chambal River have proven to negatively impact the gharial population.
Threats in unprotected stretches of the Karnali River include quarrying for boulders, sand mining and unlicensed fishing.
Conservation
The gharial is listed on CITES Appendix I. In India, it is protected under the Wildlife Protection Act of 1972. In Nepal, it is fully protected under the National Parks and Wildlife Conservation Act of 1973.
Reintroduction programmes
Since the late 1970s, the gharial conservation approach was focused on reintroduction. Rivers in protected areas in India and Nepal used to be restocked with captive bred juvenile gharials. Gharial eggs were incubated, hatched and juvenile gharials raised for two to three years and released when about one metre in length.
In 1975, the Indian Crocodile Conservation Project was set up under the auspices of the Government of India, initially in Odisha's Satkosia Gorge Sanctuary. It was implemented with financial aid of the United Nations Development Fund and the Food and Agriculture Organization. The country's first gharial breeding center was built in Nandankanan Zoological Park. A male gharial was flown in from Frankfurt Zoological Garden to become one of the founding animals of the breeding program. In subsequent years, several protected areas were established. In 1976, two breeding centres were established in Uttar Pradesh, one in Kukrail Reserve Forest and one in Katarniaghat Wildlife Sanctuary, with facilities to hatch and raise up to 800 gharials each year for release in rivers.
Between 1975 and 1982, sixteen crocodile rehabilitation centers and five crocodile sanctuaries were established in the country. Gharial eggs were initially purchased from Nepal. In 1991, the Ministry of Environment and Forests withdrew funds for the captive-breeding and egg-collection programs, arguing that the project had served its purpose. In 1997–1998, over 1,200 gharials and over 75 nests were located in the National Chambal Sanctuary, but no surveys were carried out between 1999 and 2003. Gharial eggs collected from wild and captive-breeding nests amounted to 12,000 until 2004. Eggs were incubated, and hatchlings were reared to a length of about one meter or more.
More than 5,000 gharials were released into Indian rivers between the early 1980s and 2006. Despite the release of 142 gharials between 1982 and 2007 into the Ken River, only one adult female gharial was observed in the river in spring 2013, indicating that most of the released gharials had not reproduced.
Juvenile gharials have also been released into the Beas River in Punjab, India.
In Nepal, wild eggs collected along rivers have been incubated in the Gharial Conservation and Breeding Center in Chitwan National Park since 1978. The first batch of 50 gharials was released in spring 1981 into the Narayani River. In subsequent years, gharials were also released into five other rivers in the country. In 2016, this center was overcrowded with more than 600 gharials aged between 5 and 12 years, and many were too old to be released.
Between 1981 and 2018, a total of 1,365 gharials were released in the Rapti–Narayani river system.
Reintroducing gharials helped to maintain this population, but the survival rate of released gharials was rather low. Of 36 marked gharials released in the spring seasons of 2002 and 2003 into the Rapti–Narayani rivers, only 14 were found alive in spring 2004. This reintroduction programme has been criticised in 2017 as not being comprehensive and coordinated, as often too old and unsexed gharials were released at disturbed localities during unfavourable cold months and without assessing the efficiency of these releases.
It has been suggested to instead leave wild nests in place, increase protection of nesting and basking sites and monitor the movement of gharials.
Releasing captive-reared gharials did not contribute significantly to re-establishing viable populations.
Monitoring of released gharials revealed that the reintroduction programmes did not address multiple factors affecting their survival. These factors include disturbances from diversions of river courses, sand mining, cultivation of riversides, fishing by local people and mortality related to fishing methods like the use of gill nets and dynamite.
In 2017, members of the Crocodile Specialist Group therefore recommended to foster engagement of local communities in gharial conservation programs.
In May 2023, sightings of the Gharial were reported in the Punjab region of Pakistan. This marked the first confirmed sighting of the species in Pakistan after a presumed absence of three decades. In response to these sightings, WWF-Pakistan, in collaboration with the other partners, aims to step up conservation efforts for the Gharial. The goal is to ensure that the newly discovered population not only survives but thrives. Pakistan has requested the transfer of hundreds of Gharial crocodiles from Nepal in an effort to reintroduce the species.
In situ projects
The riverbanks of Girwa river were cleared from woody vegetation on sand banks and mid-river islands in 2019, and sand was added in 2020 to create an artificial sand bank of about . This intervention helped to stabilise and optimise the soil temperature at this site. In 2020, the number of gharial nests on this river stretch increased to 36 from 25 in 2018, and the number of unhatched eggs and dead hatchlings decreased significantly.
In captivity
As of 1999, gharials were also kept in the Madras Crocodile Bank Trust, Mysore Zoo, Jaipur Zoo and Kukrail Gharial Rehabilitation Centre in India.
In Europe, gharials are kept in Prague Zoo and Protivin Crocodile Zoo in the Czech Republic, and the Berlin Zoo in Germany.
La Ferme aux Crocodiles, a crocodile farm in France, received six juveniles in 2000 from the Gharial Breeding Centre in Nepal.
In the United States, gharials are kept in Busch Gardens Tampa, Cleveland Metroparks Zoo, Fort Worth Zoo, Honolulu Zoo, San Diego Zoo, National Zoological Park, San Antonio Zoo and Aquarium and St. Augustine Alligator Farm Zoological Park. Bronx Zoo and Los Angeles Zoo received gharials in 2017. In 2023, Fort Worth Zoo announced the birth of four gharials.
In culture
The earliest known depictions of the gharial date to the Indus Valley civilisation. Seals and tablets show gharials with fish in their mouths and surrounded by fish. A tablet shows a deity flanked by a gharial and a fish. These pieces are about 4,000 years old and were found at Mohenjo-daro and Amri, Sindh.
A gharial is depicted on one of the rock carvings on a pillar of the Sanchi Stupa, which dates to the 3rd century BC.
In Hindu mythology, the gharial is the vehicle of the river deity Gaṅgā and of the wind and sea deity Varuna.
In the 16th-century book Baburnama, Zahir-ud-din Muhammad Babur accounted of a gharial sighting in the Ghaghara River between Ghazipur and Benares in 1526.
In 1915, a British officer observed the traditional method of Kehal fishermen hunting gharials along the Indus. They staked nets about below the waterline close by a sandbank and waited hidden for gharials to come out of the river for basking. After some time, they left their hiding places, prompting the gharials to dart off to the river and get entangled in the nets.
Local people in Nepal attributed various mystical powers to the ghara of male gharials and killed them to collect their snouts. Tharu people believed that the ghara would repel insects and pests when burnt in a field, and that gharial eggs would be an effective cough medicine and aphrodisiac. Jewellery found in gharial stomachs may have been the reason for the belief of local people that they would eat humans.
Local names for the gharial include 'Lamthore gohi' and 'Chimpta gohi' in Nepali, whereby gohi means crocodile; 'Gharial' in Hindi; 'Susar' in Marathi; 'Nakar' and 'Bahsoolia nakar' in Bihari; 'Thantia kumhira' in Odia, with 'thantia' being derived from the Sanskrit word 'tuṇḍa' meaning beak, snout, elephant's trunk; the male is called 'Ghadiala' and the female 'Thantiana' in Odia.
| Biology and health sciences | Crocodilia | Animals |
1347945 | https://en.wikipedia.org/wiki/Great%20Filter | Great Filter | The Great Filter is the idea that, in the development of life from the earliest stages of abiogenesis to reaching the highest levels of development on the Kardashev scale, there is a barrier to development that makes detectable extraterrestrial life exceedingly rare. The Great Filter is one possible resolution of the Fermi paradox.
The concept originates in Robin Hanson's argument that the failure to find any extraterrestrial civilizations in the observable universe implies that something is wrong with one or more of the arguments (from various scientific disciplines) that the appearance of advanced intelligent life is probable; this observation is conceptualized in terms of a "Great Filter" which acts to reduce the great number of sites where intelligent life might arise to the tiny number of intelligent species with advanced civilizations actually observed (currently just one: human). This probability threshold, which could lie in the past or following human extinction, might work as a barrier to the evolution of intelligent life, or as a high probability of self-destruction. The main conclusion of this argument is that the easier it was for life to evolve to the present stage, the bleaker the future chances of humanity probably are.
The idea was first proposed in an online essay titled "The Great Filter – Are We Almost Past It?". The first version was written in August 1996 and the article . Hanson's formulation has received recognition in several published sources discussing the Fermi paradox and its implications.
Main argument
Fermi paradox
There is no reliable evidence that aliens have visited Earth; we have observed no intelligent extraterrestrial life with current technology, nor has SETI found any transmissions from other civilizations. The Universe, apart from the Earth, seems "dead"; Hanson states:
Our planet and solar system, however, don't look substantially colonized by advanced competitive life from the stars, and neither does anything else we see. To the contrary, we have had great success at explaining the behavior of our planet and solar system, nearby stars, our galaxy, and even other galaxies, via simple "dead" physical processes, rather than the complex purposeful processes of advanced life. Life is expected to expand to fill all available niches. With technology such as self-replicating spacecraft, these niches would include neighboring star systems and even, on longer time scales which are still small compared to the age of the universe, other galaxies. Hanson notes, "If such advanced life had substantially colonized our planet, we would know it by now."
The Great Filter
With no evidence of intelligent life in places other than Earth, it appears that the process of starting with a star and ending with "advanced explosive lasting life" must be unlikely. This implies that at least one step in this process must be improbable. Hanson's list, while incomplete, describes the following nine steps in an "evolutionary path" that results in the colonization of the observable universe:
The right star system (including organics and potentially habitable planets)
Reproductive molecules (e.g. RNA)
Simple (prokaryotic) single-cell life
Complex (eukaryotic) single-cell life
Sexual reproduction
Multi-cell life
Tool-using animals with intelligence
A civilization advancing toward the potential for a colonization explosion (where we are now)
Colonization explosion
According to the Great Filter hypothesis, at least one of these steps—if the list were complete—must be improbable. If it is not an early step (i.e., in the past), then the implication is that the improbable step lies in the future and humanity's prospects of reaching step 9 (interstellar colonization) are still bleak. If the past steps are likely, then many civilizations would have developed to the current level of the human species. However, none appear to have made it to step 9, or the Milky Way would be full of colonies. So perhaps step 9 is the unlikely one, and the only things that appear likely to keep us from step 9 are some sort of catastrophe, an underestimation of the impact of procrastination as technology increasingly unburdens existence, or resource exhaustion leading to the impossibility of making the step due to consumption of the available resources (for example highly constrained energy resources). So by this argument, finding multicellular life on Mars (provided it evolved independently) would be bad news, since it would imply steps 2–6 are easy, and hence only 1, 7, 8 or 9 (or some unknown step) could be the big problem.
Although steps 1–8 have occurred on Earth, any one of these may be unlikely. If the first seven steps are necessary preconditions to calculating the likelihood (using the local environment) then an anthropically biased observer can infer nothing about the general probabilities from its (pre-determined) surroundings.
In a 2020 paper, Jacob Haqq-Misra, Ravi Kumar Kopparapu, and Edward Schwieterman argued that current and future telescopes searching for biosignatures in the ultraviolet to near-infrared wavelengths could place upper bounds on the fraction of planets in the galaxy that host life. Meanwhile, the evolution of telescopes that can detect technosignatures at mid-infrared wavelengths could provide insights into the Great Filter. They say that if planets with technosignatures are abundant, then this can increase confidence that the Great Filter is in the past. On the other hand, if finding that life is commonplace while technosignatures are absent, then this would increase the likelihood that the Great Filter lies in the future.
Recently, paleobiologist Olev Vinn has suggested that the great filter may exist between steps 8 and 9 due to inherited behavior patterns (IBP) that initially occur in all intelligent biological organisms. These IBPs are incompatible with conditions prevailing in technological civilizations and could inevitably lead to the self-destruction of civilization in multiple ways.
In a specific formulation named the "Berserker hypothesis", a filter exists between steps 8 and 9 in which each civilization is destroyed by a lethal Von Neumann probe created by a more advanced civilization.
Responses
There are many alternative scenarios that might allow for the evolution of intelligent life to occur multiple times without either catastrophic self-destruction or glaringly visible evidence. These are possible resolutions to the Fermi paradox: "They do exist, but we see no evidence". Other ideas include: it is too expensive to spread physically throughout the galaxy; Earth is purposely isolated; it is dangerous to communicate and hence civilizations actively hide, among others.
Astrobiologists Dirk Schulze-Makuch and William Bains, reviewing the history of life on Earth, including convergent evolution, concluded that transitions such as oxygenic photosynthesis, the eukaryotic cell, multicellularity, and tool-using intelligence are likely to occur on any Earth-like planet given enough time. They argue that the Great Filter may be abiogenesis, the rise of technological human-level intelligence, or an inability to settle other worlds because of self-destruction or a lack of resources.
Astronomer Seth Shostak of the SETI Institute argues that one can postulate a galaxy filled with intelligent extraterrestrial civilizations that have failed to colonize Earth. Perhaps the aliens lacked the intent and purpose to colonize or depleted their resources, or maybe the galaxy is colonized but in a heterogeneous manner, or the Earth could be located in a "galactic backwater". Although absence of evidence generally is only weak evidence of absence, the absence of extraterrestrial megascale engineering projects, for example, might point to the Great Filter at work. Does this mean that one of the steps leading to intelligent life is unlikely? According to Shostak:
This is, of course, a variant on the Fermi paradox: We don't see clues to widespread, large-scale engineering, and consequently we must conclude that we're alone. But the possibly flawed assumption here is when we say that highly visible construction projects are an inevitable outcome of intelligence. It could be that it's the engineering of the small, rather than the large, that is inevitable. This follows from the laws of inertia (smaller machines are faster, and require less energy to function) as well as the speed of light (small computers have faster internal communication). It may be—and this is, of course, speculation—that advanced societies are building small technology and have little incentive or need to rearrange the stars in their neighborhoods, for instance. They may prefer to build nanobots instead. It should also be kept in mind that, as Arthur C. Clarke said, truly advanced engineering would look like magic to us—or be unrecognizable altogether. By the way, we've only just begun to search for things like Dyson spheres, so we can't really rule them out.Joseph Voros in "Macro-Perspectives Beyond the World System" (2007) points out that some researchers have attempted to search for energy signatures that could be traced to Dyson-like structures (shells, swarms, or spheres). So far, none have been found. See for example, Tilgner & Heinrichsen, "A Program to Search for Dyson Spheres with the Infrared Space Observatory", Acta Astronautica Vol. 42 (May–June, 1998), pp. 607–612; and Timofeev et al. "A search of the IRAS database for evidence of Dyson Spheres", Acta Astronautica Vol. 46, (June 2000), pp. 655–659.
| Physical sciences | Astronomy basics | Astronomy |
1348347 | https://en.wikipedia.org/wiki/Funding%20of%20science | Funding of science | Research funding is a term generally covering any funding for scientific research, in the areas of natural science, technology, and social science. Different methods can be used to disburse funding, but the term often connotes funding obtained through a competitive process, in which potential research projects are evaluated and only the most promising receive funding. It is often measured via Gross domestic expenditure on R&D (GERD).
Most research funding comes from two major sources: corporations (through research and development departments) and government (primarily carried out through universities and specialized government agencies; often known as research councils). A smaller amount of scientific research is funded by charitable foundations, especially in relation to developing cures for diseases such as cancer, malaria, and AIDS.
According to the Organisation for Economic Co-operation and Development (OECD), more than 60% of research and development in scientific and technical fields is carried out by industry, and 20% and 10% respectively by universities and government. Comparatively, in countries with less GDP such as Portugal and Mexico, the industry contribution is significantly lower. The government funding proportion in certain industries is higher, and it dominates research in social science and humanities. In commercial research and development, all but the most research-oriented corporations focus more heavily on near-term commercialization possibilities rather than "blue-sky" ideas or technologies (such as nuclear fusion).
History
Conducting research requires funds. The funding trend for research has gone from a closed patronage system to which only few could contribute, to an open system with multiple funding possibilities.
In the early Zhou dynasty (-c. 6th century to 221 BCE), government officials used their resources to fund schools of thought of which they were patron. The bulk of their philosophies are still relevant, including Confucianism, Legalism and Taoism.
During the Mayan Empire (-c. 1200–1250), scientific research was funded for religious purposes. Research there developed a Venus Table, showing precise astronomical data about the position of Venus in the sky. In Cairo (-c. 1283), the Mamluk Sultan Qalawun funded a monumental hospital, patronizing the medical sciences over the religious sciences. Furthermore, Tycho Brahe was given an estate (-c. 1576 – 1580) by his royal patron King Frederik II, which was used to build Uraniborg, an early research institute.
The age of the academies
In 1700–1799, scientific academies became central creators of scientific knowledge. Funded by state sponsorship, academic societies were free to manage scientific developments. Membership was exclusive in terms of gender, race and class, but academies opened the world of research up beyond the traditional patronage system.
In 1799, Louis-Nicolas Robert patented the paper machine. When he quarreled over invention ownership, he sought financing from the Fourdrinier brothers. In 19th century Europe, businessmen financed the application of science to industry.
In the eighteenth and nineteenth centuries, as the pace of technological progress increased before and during the Industrial Revolution, most scientific and technological research was carried out by individual inventors using their own funds. A system of patents was developed to allow inventors a period of time (often twenty years) to commercialize their inventions and recoup a profit, although in practice many found this difficult.
The Manhattan Project (1942 – 1946) had cost $27 billion and employed 130,000 people, many of them scientists charged with producing the first nuclear weapons. In 1945, 70 scientists signed the Szilard petition, asking President Truman to make a demonstration of the power of the bomb before using it. Most of the signers lost their jobs in military research.
In the twentieth century, scientific and technological research became increasingly systematized, as corporations developed, and discovered that continuous investment in research and development could be a key element of competitive success. It remained the case, however, that imitation by competitors - circumventing or simply flouting patents, especially those registered abroad - was often just as successful a strategy for companies focused on innovation in matters of organization and production technique, or even in marketing.
In 2025, many funders make research outcomes transparent and accessible in data repositories or Open-access. Some researchers turn to crowdfunding in search of new projects to fund. Private and public foundations, governments, and others sponsor opportunities for researchers. As new funding sources become available, the research community grows and becomes accessible to a wider, and more diverse group of scientists.
Methodology to measure science funding
The guidelines for R&D data collections are laid down in the Frascati Manual published by the OECD. In the publication, R&D denotes three type of activity: basic research, applied research and experimental development. This definition does not cover innovation but it may feed into the innovative process. Business sector innovation has a dedicated OECD manual.
The most frequently used measurement for R&D is Gross domestic expenditure on R&D (GERD). GERD is often represented in GERD-to-GDP ratios, as it allows for easier comparisons between countries. The data collection for GERD is based on reporting by performers. GERD differentiates according to the funding sector (business, enterprise, government, higher education, private non-profit, rest of the world) and the sector of performance (all funding sectors with the exception of rest of the world as GERD only measures activity within the territory of a country). The two may coincide for example when government funds government performed R&D.
Government funded science also may be measured by the Government budget appropriations and outlays for R&D (GBAORD/ GBARD). GBARD is a funder-based method, it denotes what governments committed to R&D (even if final payment might be different). GERD-source of funding-government and GBARD are not directly comparable. On data collection, GERD is performer based, GBARD is funder. The level of government considered also differs: GERD may include spending by all levels of the government (federal – state – local), whereas GBARD excludes the local level and often lacks state level data. On geographic coverage, GERD takes into account performance within the territory of a country whereas GBARD also payments to the Rest of the world.
Comparisons on the effectiveness of both the different sources of funding and sectors of performance as well as their interplay have been made. The analysis often boils down to whether public and private finance show crowding-in or crowding-out patterns.
Funding types: public and private
Public/State Funding
Public funding refers to activities financed by tax-payers money. This is primarily the case when the source of funds is channeled through government agencies. Higher education institutions are usually not completely publicly financed as they charge tuition fees and may receive funds from non-public sources.
Rationale for funding
R&D is a costly, and long-term investment to which disruptions are harmful.
The public sector has multiple reasons to fund science. The private sector is said to focus on the closer to the market stage of R&D policy, where appropriability hence private returns are high. Basic research is weak on appropriability and so remains risky and under-financed. Consequently, although governmental sponsorship of research may provide support across the R&D value chain, it is often characterized as a market failure induced intervention. Market incentives to invest in early-stage research are low. The theory of public goods seconds this argument. Publicly funded research often supports research fields where social rate of return may be higher than private rate of return. Appropriability potential is the potential for an entity to capture the value of an innovation or research outcome. The general free rider problem of public goods is a threat especially in case of global public goods such as climate change research, which may lower incentives to invest by both the private sector but also other governments.
In endogenous growth theories, R&D contributes to growth. Some have depicted this relationship in the inverse, claiming that growth drives innovation. As of 2013, science workers applying their (tacit) knowledge may be considered an economic driver. When this knowledge and/or human capital emigrates, countries face the so-called brain–drain. Science policy can assist to avoid this as large shares of governmental R&D is spent on researchers and supporting staff personnel salaries. In this sense, science funding is not only discretionary spending but also has elements of entitlement spending.
R&D funded and especially performed by the State may allow greater influence over its direction. This is particularly important in the case of R&D contributing to public goods. However, the ability of governments have been criticized over whether they are best positioned to pick winners and losers. In the EU, dedicated safeguards have been enacted under a dedicated form of competition law called State Aid. State Aid safeguards business activities from governmental interventions. This invention was largely driven by the German ordoliberal school as to eliminate state subsidies advocated by the French dirigiste. Threats to global public goods has refueled the debate on the role of governments beyond a mere market failure fixer, the so-called mission-driven policies.
Funding modalities
Governments may fund science through different instruments such as: direct subsidies, tax credits, loans, financial instruments, regulatory measures, public procurement etc.
While direct subsidies have been the prominent instrument to fund business R&D, since the financial crisis a shift has taken place in OECD countries in the direction of tax breaks. The explanation seems to lay in the theoretical argument that firms know better, and in the practical benefit of lower administrative burden of such schemes. Depending on the funding type, different modalities to distribute the research funds may be used. For regulatory measures, often the competition/antitrust authorities will rule on exemptions. In case of block funding the funds may be directly allocated to given institutions such as higher education institutions with relative autonomy over their use. For competitive grants, governments are often assisted by research councils to distribute the funds. Research councils are (usually public) bodies that provide research funding in the form of research grants or scholarships. These include arts councils and research councils for the funding of science.
List of research councils
An incomplete list of national and international pan-disciplinary public research councils:
Conditionality
In addition to project deliverables, funders also increasingly introduce new eligibility requirements alongside traditional ones such as research integrity/ethics.
The 2016 Open Science movement, tied funding increasingly tied to data management plans and making data FAIR. The Open Science requirement complements Open Access mandates which in 2025 are widespread.
The gender dimension also gained ground in recent years. The European Commission mandates research applicants to adopt gender equality plans across their organization. The UK Research and Innovation Global Challenges Research Fund mandates a gender equality statement.
As of 2022, the European Commission also introduced a “Do No Significant Harm” principle to the Framework Program which aims to curb the environmental footprint of scientific projects. "Do No Significant Harm" has been criticized as coupled with other eligibility requirements it is often characterized as red-tape. Since 2020, European Commission has been trying to simplify the Framework Program with limited success. Simplification attempts were also taken by the UK Research and Innovation.
Process
Often scientists apply for research funding which a granting agency may (or may not) approve to financially support. These grants require a lengthy process as the granting agency can inquire about the researcher(s)'s background, the facilities used, the equipment needed, the time involved, and the overall potential of the scientific outcome. The process of grant writing and grant proposing is a somewhat delicate process for both the grantor and the grantee: the grantors want to choose the research that best fits their scientific principles, and the individual grantees want to apply for research in which they have the best chances but also in which they can build a body of work towards future scientific endeavors.
As of 2009, the Engineering and Physical Sciences Research Council in the United Kingdom devised an alternative method of fund-distribution: the sandpit.
Most universities have research administration offices to facilitate the interaction between the researcher and the granting agency. "Research administration is all about service—service to our faculty, to our academic units, to the institution, and to our sponsors. To be of service, we first have to know what our customers want and then determine whether or not we are meeting those needs and expectations."
In the United States of America, the National Council of University Research Administrators serves its members and advances the field of research administration through education and professional development programs, the sharing of knowledge and experience, and by fostering a professional, collegial, and respected community.
Hard money versus soft money
In academic contexts, hard money may refer to funding received from a government or other entity at regular intervals, thus providing a steady inflow of financial resources to the beneficiary. The antonym, soft money, refers to funding provided only through competitive research grants and the writing of grant proposals.
Hard money is usually issued by the government for the advancement of certain projects or for the benefit of specific agencies. Community healthcare, for instance, may be supported by the government by providing hard money. Since funds are disbursed regularly and continuously, the offices in charge of such projects are able to achieve their objectives more effectively than if they had been issued one-time grants.
Individual jobs at a research institute may be classified as "hard-money positions" or "soft-money positions"; the former are expected to provide job security because their funding is secure in the long term, whereas individual "soft-money" positions may come and go with fluctuations in the number of grants awarded to an institution.
Private funding: industrial/philanthropy/crowdfunding
Private funding for research comes from philanthropists, crowd-funding, private companies, non-profit foundations, and professional organizations. Philanthropists and foundations have been pouring millions of dollars into a wide variety of scientific investigations, including basic research discovery, disease cures, particle physics, astronomy, marine science, and the environment. Privately funded research has been adept at identifying important and transformative areas of scientific research. Many large technology companies spend billions of dollars on research and development each year to gain an innovative advantage over their competitors, though only about 42% of this funding goes towards projects that are considered substantially new, or capable of yielding radical breakthroughs. New scientific start-up companies initially seek funding from crowd-funding organizations, venture capitalists, and angel investors, gathering preliminary results using rented facilities, but aim to eventually become self-sufficient.
Europe and the United States have both reiterated the need for further private funding within universities. The European Commission highlights the need for private funding via research in policy areas such the European Green Deal and Europe's role in the digital age.
Influence on research
The source of funding may introduce conscious or unconscious biases into a researcher's work. This is highly problematic due to academic freedom in case of universities and regulatory capture in case of government-funded R&D.
Conflict of Interest
Disclosure of potential conflicts of interest (COIs) is used by journals to guarantee credibility and transparency of the scientific process. Conflict of interest disclosure, however, is not systematically nor consistently dealt with by journals that publish scientific research results.
When research is funded by the same agency that can be expected to gain from a favorable outcome there is a potential for biased results and research shows that results are indeed more favorable than would be expected from a more objective view of the evidence. A 2003 systematic review studied the scope and impact of industry sponsorship in biomedical research. The researchers found financial relationships among industry, scientific investigators, and academic institutions widespread. Results showed a statistically significant association between industry sponsorship and pro-industry conclusions and concluded that "Conflicts of interest arising from these ties can influence biomedical research in important ways". A British study found that a majority of the members on national and food policy committees receive funding from food companies.
In an effort to cut costs, the pharmaceutical industry has turned to the use of private, nonacademic research groups (i.e., contract research organizations [CROs]) which can do the work for less money than academic investigators. In 2001 CROs came under criticism when the editors of 12 major scientific journals issued a joint editorial, published in each journal, on the control over clinical trials exerted by sponsors, particularly targeting the use of contracts which allow sponsors to review the studies prior to publication and withhold publication of any studies in which their product did poorly. They further criticized the trial methodology stating that researchers are frequently restricted from contributing to the trial design, accessing the raw data, and interpreting the results.
The Cochrane Collaboration, a worldwide group that aims to provide compiled scientific evidence to aid well informed health care decisions, conducts systematic reviews of randomized controlled trials of health care interventions and tries to disseminate the results and conclusions derived from them. A few more recent reviews have also studied the results of non-randomized, observational studies. The systematic reviews are published in the Cochrane Library. A 2011 study done to disclose possible conflicts of interests in underlying research studies used for medical meta-analyses reviewed 29 meta-analyses and found that conflicts of interest in the studies underlying the meta-analyses were rarely disclosed. The 29 meta-analyses reviewed an aggregate of 509 randomized controlled trials. Of these, 318 trials reported funding sources with 219 (69%) industry funded. 132 of the 509 trials reported author disclosures of conflict of interest, with 91 studies (69%) disclosing industry financial ties with one or more authors. However, the information was seldom reflected in the meta-analyses. Only two (7%) reported funding sources and none reported author-industry ties. The authors concluded, "without acknowledgment of COI due to industry funding or author industry financial ties from RCTs included in meta-analyses, readers' understanding and appraisal of the evidence from the meta-analysis may be compromised."
In 2003 researchers looked at the association between authors' published positions on the safety and efficacy in assisting with weight loss of olestra, a fat substitute manufactured by the Procter & Gamble (P&G), and their financial relationships with the food and beverage industry. They found that supportive authors were significantly more likely than critical or neutral authors to have financial relationships with P&G and all authors disclosing an affiliation with P&G were supportive. The authors of the study concluded: "Because authors' published opinions were associated with their financial relationships, obtaining noncommercial funding may be more essential to maintaining objectivity than disclosing personal financial interests."
A 2005 study in the journal Nature surveyed 3247 US researchers who were all publicly funded (by the National Institutes of Health). Out of the scientists questioned, 15.5% admitted to altering design, methodology or results of their studies due to pressure of an external funding source.
Regulatory capture
Private funding also may be channelled to public funders. In 2022, a news story broke following the resignation of Eric Lander, former director of the Office of Science and Technology Policy (OSTP) at the Biden administration, that the charity of former Google executive Eric Schmidt, Schmidt Futures, paid the salary of a number employees of the OSTP. Ethics inquiries were initiated in the OSTP.
Efficiency of funding
The traditional measurement for efficiency of funding are publication output, citation impact, number of patents, number of PhDs awarded etc. However, the use of journal impact factor has generated a publish-or-perish culture and a theoretical model has been established whose simulations imply that peer review and over-competitive research funding foster mainstream opinion to monopoly. Calls have been made to reform research assessment, most notably in the San Francisco Declaration on Research Assessment and the Leiden Manifesto for research metrics. The current system also has limitations to measure excellence in the Global South. Novel measurement systems such as the Research Quality Plus has been put forward to better emphasize local knowledge and contextualization in the evaluation of excellence.
Another question is how to allocate funds to different disciplines, institutions, or researchers. A recent study by Wayne Walsh found that “prestigious institutions had on average 65% higher grant application success rates and 50% larger award sizes, whereas less-prestigious institutions produced 65% more publications and had a 35% higher citation impact per dollar of funding.”
Trends
In endogenous growth theories R&D contributes to economic growth. Therefore, countries have strong incentives to maintain investments in R&D.
By country
Different countries spend vastly different amounts on research, in both absolute and relative terms. For instance, South Korea and Israel spend more than 4% of their GDP while many less developed countries spend less than 1%. In developed economies, GERD is financed mainly by the business sector, whereas the government and the university sector dominates in less-developed economies. In some countries, funding from the Rest of the World makes up 20-30% of total GERD, probably due to FDI and foreign aid, but only in Mali it is the main source of fund. Private non-profit is not the main source of fund in any countries, but it reaches 10% of total GERD in Columbia and Honduras.
When comparing annual GERD and GDP Growth, it can be seen that countries with lower GERD are often growing faster. However, as most of these countries are developing, their growth is probably driven by other factors of production. On the other hand, developed countries who have higher GERD also produce positive growth rates. GERD in these countries has a more substantial contribution to growth rate.
Recessions
In crisis, business R&D tends to act procyclically. As R&D is a long-term investments and so disruptions should be avoided Keynesian countercyclical reactions were advocated for in the wake of the 2008 financial crisis, but this was difficult to achieve for some countries. Due to the nature of COVID-19, the pandemic accelerated publicly funded R&D spending in 2020, primarily into the pharmaceutical industry. A fall is expected in spending for 2021, although not below 2020 levels. The pandemic made health research and sectors with strategic value-chain dependencies the main target of science funding.
| Physical sciences | Science basics | Basics and measurement |
1348798 | https://en.wikipedia.org/wiki/Skolem%27s%20paradox | Skolem's paradox | In mathematical logic and philosophy, Skolem's paradox is the apparent contradiction that a countable model of first-order set theory could contain an uncountable set. The paradox arises from part of the Löwenheim–Skolem theorem; Thoralf Skolem was the first to discuss the seemingly contradictory aspects of the theorem, and to discover the relativity of set-theoretic notions now known as non-absoluteness. Although it is not an actual antinomy like Russell's paradox, the result is typically called a paradox and was described as a "paradoxical state of affairs" by Skolem.
In model theory, a model corresponds to a specific interpretation of a formal language or theory. It consists of a domain (a set of objects) and an interpretation of the symbols and formulas in the language, such that the axioms of the theory are satisfied within this structure. The Löwenheim–Skolem theorem shows that any model of set theory in first-order logic, if it is consistent, has an equivalent model that is countable. This appears contradictory, because Georg Cantor proved that there exist sets which are not countable. Thus the seeming contradiction is that a model that is itself countable, and which therefore contains only countable sets, satisfies the first-order sentence that intuitively states "there are uncountable sets".
A mathematical explanation of the paradox, showing that it is not a true contradiction in mathematics, was first given in 1922 by Skolem. He explained that the countability of a set is not absolute, but relative to the model in which the cardinality is measured. Skolem's work was harshly received by Ernst Zermelo, who argued against the limitations of first-order logic and Skolem's notion of "relativity," but the result quickly came to be accepted by the mathematical community.
The philosophical implications of Skolem's paradox have received much study. One line of inquiry questions whether it is accurate to claim that any first-order sentence actually states "there are uncountable sets". This line of thought can be extended to question whether any set is uncountable in an absolute sense. More recently, scholars such as Hilary Putnam have introduced the paradox and Skolem's concept of relativity to the study of the philosophy of language.
Background
One of the earliest results in set theory, published by Cantor in 1874, was the existence of different sizes, or cardinalities, of infinite sets. An infinite set is called countable if there is a function that gives a one-to-one correspondence between and the natural numbers, and is uncountable if there is no such correspondence function. In 1874, Cantor proved that the real numbers were uncountable; in 1891, he proved by his diagonal argument the more general result known as Cantor's theorem: for every set , the power set of cannot be in bijection with itself. When Zermelo proposed his axioms for set theory in 1908, he proved Cantor's theorem from them to demonstrate their strength.
In 1915, Leopold Löwenheim gave the first proof of what Skolem would prove more generally in 1920 and 1922, the Löwenheim–Skolem theorem. Löwenheim showed that any first-order sentence with a model also has a model with a countable domain; Skolem generalized this to infinite sets of sentences. The downward form of the Löwenheim–Skolem theorem shows that if a countable first-order collection of axioms is satisfied by an infinite structure, then the same axioms are satisfied by some countably infinite structure. Since the first-order versions of standard axioms of set theory (such as Zermelo–Fraenkel set theory) are a countable collection of axioms, this implies that if these axioms are satisfiable, they are satisfiable in some countable model.
The result and its implications
In 1922, Skolem pointed out the seeming contradiction between the Löwenheim–Skolem theorem, which implies that there is a countable model of Zermelo's axioms, and Cantor's theorem, which states that uncountable sets exist, and which is provable from Zermelo's axioms. "So far as I know," Skolem wrote, "no one has called attention to this peculiar and apparently paradoxical state of affairs. By virtue of the axioms we can prove the existence of higher cardinalities... How can it be, then, that the entire domain B [a countable model of Zermelo's axioms] can already be enumerated by means of the finite positive integers?"
However, this is only an apparent paradox. In the context of a specific model of set theory, the term "set" does not refer to an arbitrary set, but only to a set that is actually included in the model. The definition of countability requires that a certain one-to-one correspondence between a set and the natural numbers must exist. This correspondence itself is a set. Skolem resolved the paradox by concluding that such a set does not necessarily exist in a countable model; that is, countability is "relative" to a model, and countable, first-order models are incomplete.
Though Skolem gave his result with respect to Zermelo's axioms, it holds for any standard first-order theory of sets, such as ZFC. Consider Cantor's theorem as a long formula in the formal language of ZFC. If ZFC has a model, call this model and its domain . The interpretation of the element symbol , or , is a set of ordered pairs of elements of in other words, is a subset of . Since the Löwenheim–Skolem theorem guarantees that is countable, then so must be . Two special elements of model the natural numbers and the power set of the natural numbers . There is only a countably infinite set of ordered pairs in of the form , because is countable. That is, only countably many elements of model members of the uncountable set . However, there is no contradiction with Cantor's theorem, because what it states is simply that no element of models a bijective function from to .
Skolem used the term "relative" to describe when the same set could be countable in one model of set theory and not countable in another: relative to one model, no enumerating function can put some set into correspondence with the natural numbers, but relative to another model, this correspondence may exist. He described this as the "most important" result in his 1922 paper. Contemporary set theorists describe concepts that do not depend on the choice of a transitive model as absolute. From their point of view, Skolem's paradox simply shows that countability is not an absolute property in first-order logic.
Skolem described his work as a critique of (first-order) set theory, intended to illustrate its weakness as a foundational system:
Reception by the mathematical community
It took some time for the theory of first-order logic to be developed enough for mathematicians to understand the cause of Skolem's result; no resolution of the paradox was widely accepted during the 1920s. In 1928, Abraham Fraenkel still described the result as an antinomy:
In 1925, John von Neumann presented a novel axiomatization of set theory, which developed into NBG set theory. Very much aware of Skolem's 1922 paper, von Neumann investigated countable models of his axioms in detail. In his concluding remarks, von Neumann commented that there is no categorical axiomatization of set theory, or any other theory with an infinite model. Speaking of the impact of Skolem's paradox, he wrote:
Zermelo at first considered Skolem's paradox a hoax, and he spoke against Skolem's "relativism" in 1931. Skolem's result applies only to what is now called first-order logic, but Zermelo argued against the finitary metamathematics that underlie first-order logic, as Zermelo was a mathematical Platonist who opposed intuitionism and finitism in mathematics. Zermelo believed in a kind of infinite Platonic ideal of logic, and he held that mathematics had an inherently infinite character. Zermelo argued that his axioms should instead be studied in second-order logic, a setting in which Skolem's result does not apply. Zermelo published a second-order axiomatization of set theory in 1930. Zermelo's further work on the foundations of set theory after Skolem's paper led to his discovery of the cumulative hierarchy and formalization of infinitary logic.
The surprise with which set theorists met Skolem's paradox in the 1920s was a product of their times. Gödel's completeness theorem and the compactness theorem, theorems which illuminate the way that first-order logic behaves and established its finitary nature, were not first proved until 1929. Leon Henkin's proof of the completeness theorem, which is now a standard technique for constructing countable models of a consistent first-order theory, was not presented until 1947. Thus, in the 1920s, the particular properties of first-order logic that permit Skolem's paradox were not yet understood. It is now known that Skolem's paradox is unique to first-order logic; if set theory is studied using higher-order logic with full semantics, then it does not have any countable models. By the time that Zermelo was writing his final refutation of the paradox in 1937, the community of logicians and set theorists had largely accepted the incompleteness of first-order logic. Zermelo left this refutation unfinished.
Later opinions
Later mathematical logicians did not view Skolem's paradox a fatal flaw in set theory. Stephen Cole Kleene described the result as "not a paradox in the sense of outright contradiction, but rather a kind of anomaly". After surveying Skolem's argument that the result is not contradictory, Kleene concluded: "there is no absolute notion of countability". Geoffrey Hunter described the contradiction as "hardly even a paradox". Fraenkel et al. claimed that contemporary mathematicians are no more bothered by the lack of categoricity of first-order theories than they are bothered by the conclusion of Gödel's incompleteness theorem: that no consistent, effective, and sufficiently strong set of first-order axioms is complete.
Other mathematicians such as Reuben Goodstein and Hao Wang have gone so far as to adopt what is called a "Skolemite" view: that not only does the Löwenheim-Skolem theorem prove that set-theoretic notions of countability are relative to a model, but that every set is countable from some "absolute" perspective. L. E. J. Brouwer was another early adherent to the idea of absolute countability, arguing from the vantage of mathematical intuitionism that all sets are countable. Both the Skolemite view and Brouwer's intuitionism stand in opposition to mathematical Platonism, but Carl Posy denies the idea that Brouwer's position was a reaction to earlier set-theoretic paradoxes. Skolem was another mathematical intuitionist, but he denied that his ideas were inspired by Brouwer.
Countable models of Zermelo–Fraenkel set theory have become common tools in the study of set theory. Paul Cohen's method for extending set theory, forcing, is often explained in terms of countable models, and was described by Akihiro Kanamori as a kind of extension of Skolem's paradox. The fact that these countable models of Zermelo–Fraenkel set theory still satisfy the theorem that there are uncountable sets is not considered a pathology; Jean van Heijenoort described it as "not a paradox...[but] a novel and unexpected feature of formal systems".
Hilary Putnam considered Skolem's result a paradox, but one of the philosophy of language rather than of set theory or formal logic. He extended Skolem's paradox to argue that not only are set-theoretic notions of membership relative, but semantic notions of language are relative: there is no "absolute" model for terms and predicates in language. Timothy Bays argued that Putnam's argument applies the downward Löwenheim-Skolem theorem incorrectly, while Tim Button argued that Putnam's claim stands despite the use or misuse of the Löwenheim-Skolem theorem. Appeals to Skolem's paradox have been made several times in the philosophy of science, with scholars making use of Skolem's idea of the relativity of model structures.
| Mathematics | Model theory | null |
2746930 | https://en.wikipedia.org/wiki/Equilibrium%20thermodynamics | Equilibrium thermodynamics | Equilibrium Thermodynamics is the systematic study of transformations of matter and energy in systems in terms of a concept called thermodynamic equilibrium. The word equilibrium implies a state of balance. Equilibrium thermodynamics, in origins, derives from analysis of the Carnot cycle. Here, typically a system, as cylinder of gas, initially in its own state of internal thermodynamic equilibrium, is set out of balance via heat input from a combustion reaction. Then, through a series of steps, as the system settles into its final equilibrium state, work is extracted.
In an equilibrium state the potentials, or driving forces, within the system, are in exact balance. A central aim in equilibrium thermodynamics is: given a system in a well-defined initial state of thermodynamic equilibrium, subject to accurately specified constraints, to calculate, when the constraints are changed by an externally imposed intervention, what the state of the system will be once it has reached a new equilibrium. An equilibrium state is mathematically ascertained by seeking the extrema of a thermodynamic potential function, whose nature depends on the constraints imposed on the system. For example, a chemical reaction at constant temperature and pressure will reach equilibrium at a minimum of its components' Gibbs free energy and a maximum of their entropy.
Equilibrium thermodynamics differs from non-equilibrium thermodynamics, in that, with the latter, the state of the system under investigation will typically not be uniform but will vary locally in those as energy, entropy, and temperature distributions as gradients are imposed by dissipative thermodynamic fluxes. In equilibrium thermodynamics, by contrast, the state of the system will be considered uniform throughout, defined macroscopically by such quantities as temperature, pressure, or volume. Systems are studied in terms of change from one equilibrium state to another; such a change is called a thermodynamic process.
Ruppeiner geometry is a type of information geometry used to study thermodynamics. It claims that thermodynamic systems can be represented by Riemannian geometry, and that statistical properties can be derived from the model. This geometrical model is based on the idea that there exist equilibrium states which can be represented by points on two-dimensional surface and the distance between these equilibrium states is related to the fluctuation between them.
| Physical sciences | Thermodynamics | Physics |
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