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21438200 | https://en.wikipedia.org/wiki/Synesthesia | Synesthesia | Synesthesia (American English) or synaesthesia (British English) is a perceptual phenomenon in which stimulation of one sensory or cognitive pathway leads to involuntary experiences in a second sensory or cognitive pathway. People with synesthesia may experience colors when listening to music, see shapes when smelling certain scents, or perceive tastes when looking at words. People who report a lifelong history of such experiences are known as synesthetes. Awareness of synesthetic perceptions varies from person to person with the perception of synesthesia differing based on an individual's unique life experiences and the specific type of synesthesia that they have. In one common form of synesthesia, known as grapheme–color synesthesia or color–graphemic synesthesia, letters or numbers are perceived as inherently colored. In spatial-sequence, or number form synesthesia, numbers, months of the year, or days of the week elicit precise locations in space (e.g., 1980 may be "farther away" than 1990), or may appear as a three-dimensional map (clockwise or counterclockwise). Synesthetic associations can occur in any combination and any number of senses or cognitive pathways.
Little is known about how synesthesia develops. It has been suggested that synesthesia develops during childhood when children are intensively engaged with abstract concepts for the first time. This hypothesis—referred to as semantic vacuum hypothesis—could explain why the most common forms of synesthesia are grapheme-color, spatial sequence, and number form. These are usually the first abstract concepts that educational systems require children to learn.
The earliest recorded case of synesthesia is attributed to the Oxford University academic and philosopher John Locke, who, in 1690, made a report about a blind man who said he experienced the color scarlet when he heard the sound of a trumpet. However, there is disagreement as to whether Locke described an actual instance of synesthesia or was using a metaphor. The first medical account came from German physician Georg Tobias Ludwig Sachs in 1812. The term is from Ancient Greek 'together' and 'sensation'.
Types
There are two overall forms of synesthesia:
projective synesthesia: seeing colors, forms, or shapes when stimulated (the widely understood version of synesthesia)
associative synesthesia: feeling a very strong and involuntary connection between the stimulus and the sense that it triggers
For example, in chromesthesia (sound to color), a projector may hear a trumpet, and see an orange triangle in space, while an associator might hear a trumpet, and think very strongly that it sounds "orange".
Synesthesia can occur between nearly any two senses or perceptual modes, and at least one synesthete, Solomon Shereshevsky, experienced synesthesia that linked all five senses. Types of synesthesia are indicated by using the notation , where x is the "inducer" or trigger experience, and y is the "concurrent" or additional experience. For example, perceiving letters and numbers (collectively called graphemes) as colored would be indicated as grapheme-color synesthesia. Similarly, when synesthetes see colors and movement as a result of hearing musical tones, it would be indicated as tone → (color, movement) synesthesia.
While nearly every logically possible combination of experiences can occur, several types are more common than others.
Auditory–tactile synesthesia
In auditory–tactile synesthesia, certain sounds can induce sensations in parts of the body. For example, someone with auditory–tactile synesthesia may experience that hearing a specific word or sound feels like touch in one specific part of the body or may experience that certain sounds can create a sensation in the skin without being touched (not to be confused with the milder general reaction known as frisson, which affects approximately 50% of the population). It is one of the least common forms of synesthesia.
Chromesthesia
Another common form of synesthesia is the association of sounds with colors. For some, everyday sounds can trigger seeing colors. For others, colors are triggered when musical notes or keys are being played. People with synesthesia related to music may also have perfect pitch because their ability to see and hear colors aids them in identifying notes or keys.
The colors triggered by certain sounds, and any other synesthetic visual experiences, are referred to as photisms.
According to Richard Cytowic, chromesthesia is "something like fireworks": voice, music, and assorted environmental sounds such as clattering dishes or dog barks trigger color and firework shapes that arise, move around, and then fade when the sound ends. Sound often changes the perceived hue, brightness, scintillation, and directional movement. Some individuals see music on a "screen" in front of their faces. For Deni Simon, music produces waving lines "like oscilloscope configurations lines moving in color, often metallic with height, width, and, most importantly, depth. My favorite music has lines that extend horizontally beyond the 'screen' area."
Individuals rarely agree on what color a given sound is. Composers Franz Liszt and Nikolai Rimsky-Korsakov famously disagreed on the colors of musical keys.
Grapheme–color synesthesia
In one of the most common forms of synesthesia, individual letters of the alphabet and numbers (collectively referred to as "graphemes") are "shaded" or "tinged" with a color. While different individuals usually do not report the same colors for all letters and numbers, studies with large numbers of synesthetes find some commonalities across letters (e.g., A is likely to be red).
Some authors have argued that the term synaesthesia may not be correct when applied to the so-called grapheme-colour synesthesia and similar phenomena in which the inducer is conceptual (e.g. a letter or number) rather than sensory (e.g. sound or color). They have postulated that the term ideasthesia is a more accurate description.
Kinesthetic synesthesia
Kinesthetic synesthesia is one of the rarest documented forms of synesthesia. This form of synesthesia is a combination of different types of synesthesia. Features appear similar to auditory–tactile synesthesia but sensations are not isolated to individual numbers or letters but complex systems of relationships.
Lexical–gustatory synesthesia
This is another form of synesthesia where certain tastes are experienced when hearing words. For example, the word basketball might taste like waffles. The documentary 'Derek Tastes of Earwax' gets its name from this phenomenon, referring to pub owner James Wannerton who experiences this particular sensation whenever he hears the name spoken. It is estimated that 0.2% of the synesthesia population has this form of synesthesia, making it one of the rarest forms.
Mirror-touch synesthesia
This is a form of synesthesia where individuals feel the same/similar sensation as another person (such as touch). For instance, when such a synesthete observes someone being tapped on their shoulder, the synesthete involuntarily feels a tap on their own shoulder as well. People with this type of synesthesia have been shown to have higher empathy levels compared to the general population. This may be related to the so-called mirror neurons present in the motor areas of the brain, which have also been linked to empathy.
Misophonia
Misophonia is a neurological disorder in which negative experiences (anger, fright, hatred, disgust) are triggered by specific sounds. Cytowic suggests that misophonia is related to, or perhaps a variety of, synesthesia. Edelstein and her colleagues have compared misophonia to synesthesia in terms of connectivity between different brain regions as well as specific symptoms. They hypothesize that "a pathological distortion of connections between the auditory cortex and limbic structures could cause a form of sound-emotion synesthesia." Studies suggest that individuals with misophonia have a normal hearing sensitivity level, but their limbic system and autonomic nervous system are constantly in a "heightened state of arousal" in which abnormal reactions to sounds will be more prevalent.
Newer studies suggest that, depending on its severity, misophonia could be associated with lower cognitive control when individuals are exposed to certain associations and triggers.
It is unclear what causes misophonia. Some scientists believe the condition could be genetic, while others believe it to be present with additional conditions. There are no current treatments for the condition, but management of symptoms involves numerous coping strategies. These strategies include avoidance of situations that could trigger the reaction, mimicking the sounds, and cancelling out the sounds by using earplugs or music. Most misophoniacs use these to "overwrite" these sounds produced by others.
Number form
A number form is a mental map of numbers that automatically and involuntarily appear whenever someone who experiences number-forms synesthesia thinks of numbers. These numbers might appear in different locations and the mapping changes and varies between individuals. Number forms were first documented and named in 1881 by Francis Galton in "The Visions of Sane Persons".
Ordinal linguistic personification
Ordinal-linguistic personification (OLP, or personification) is a form of synesthesia in which ordered sequences, such as ordinal numbers, week-day names, months, and alphabetical letters are associated with personalities or genders (). Although this form of synesthesia was documented as early as the 1890s, researchers have, until recently, paid little attention to it (see History of synesthesia research). This form of synesthesia was named "OLP" in the contemporary literature by Julia Simner and colleagues although it is now also widely recognized by the term "sequence-personality" synesthesia. Ordinal linguistic personification normally co-occurs with other forms of synesthesia such as grapheme–color synesthesia.
Spatial sequence synesthesia
Those with spatial sequence synesthesia (SSS) tend to see ordinal sequences as points in space. People with SSS may have superior memories; in one study, they were able to recall past events and memories far better and in far greater detail than those without the condition. They can also see months or dates in the space around them, but most synesthetes "see" these sequences in their mind's eye. Some people see time like a clock above and around them.
Ticker-tape synesthesia
Those with ticker-tape synesthesia mentally see written words when they are heard, sometimes on imaginary strips of paper. It has been suggested that the name "subtitled synesthesia" would better describe the phenomenon.
Other forms
Other forms of synesthesia have been reported, but little has been done to analyze them scientifically. There are at least 80 types of synesthesia.
In August 2017 a research article in the journal Social Neuroscience reviewed studies with fMRI to determine if persons who experience autonomous sensory meridian response are experiencing a form of synesthesia. While a determination has not yet been made, there is anecdotal evidence that this may be the case, based on significant and consistent differences from the control group, in terms of functional connectivity within neural pathways. It is unclear whether this will lead to ASMR being included as a form of existing synesthesia, or if a new type will be considered.
Signs and symptoms
Some synesthetes often report that they were unaware their experiences were unusual until they realized other people did not have them, while others report feeling as if they had been keeping a secret their entire lives. The automatic and ineffable nature of a synesthetic experience means that the pairing may not seem out of the ordinary. This involuntary and consistent nature helps define synesthesia as a real experience. Most synesthetes report that their experiences are pleasant or neutral, although, in rare cases, synesthetes report that their experiences can lead to a degree of sensory overload.
Though often stereotyped in the popular media as a medical condition or neurological aberration, many synesthetes themselves do not perceive their synesthetic experiences as a handicap. On the contrary, some report it as a gift an additional "hidden" sense something they would not want to miss. Most synesthetes become aware of their distinctive mode of perception in their childhood. Some have learned how to apply their ability in daily life and work. Synesthetes have used their abilities in memorization of names and telephone numbers, mental arithmetic, and more complex creative activities like producing visual art, music, and theater.
Despite the commonalities which permit the definition of the broad phenomenon of synesthesia, individual experiences vary in numerous ways. This variability was first noticed early in synesthesia research. Some synesthetes report that vowels are more strongly colored, while for others consonants are more strongly colored. Self-reports, interviews, and autobiographical notes by synesthetes demonstrate a great degree of variety in types of synesthesia, the intensity of synesthetic perceptions, awareness of the perceptual discrepancies between synesthetes and non-synesthetes, and the ways synesthesia is used in work, creative processes, and daily life.
Synesthetes are very likely to participate in creative activities. It has been suggested that individual development of perceptual and cognitive skills, in addition to one's cultural environment, produces the variety in awareness and practical use of synesthetic phenomena. Synesthesia may also give a memory advantage. In one study, conducted by Julia Simner of the University of Edinburgh, it was found that spatial sequence synesthetes have a built-in and automatic mnemonic reference. Whereas a non-synesthete will need to create a mnemonic device to remember a sequence (like dates in a diary), a synesthete can simply reference their spatial visualizations.
Mechanism
As of 2015, the neurological correlates of synesthesia had not been established.
Dedicated regions of the brain are specialized for given functions. Increased cross-talk between regions specialized for different functions may account for the many types of synesthesia. For example, the additive experience of seeing color when looking at graphemes might be due to cross-activation of the grapheme-recognition area and the color area called V4 (see figure). This is supported by the fact that grapheme–color synesthetes can identify the color of a grapheme in their peripheral vision even when they cannot consciously identify the shape of the grapheme.
An alternative possibility is disinhibited feedback or a reduction in the amount of inhibition along normally existing feedback pathways. Normally, excitation and inhibition are balanced. However, if normal feedback was not inhibited as usual, then signals feeding back from late stages of multi-sensory processing might influence earlier stages such that tones could activate vision. Cytowic and Eagleman find support for the disinhibition idea in the so-called acquired forms of synesthesia that occur in non-synesthetes under certain conditions: temporal lobe epilepsy, head trauma, stroke, and brain tumors. They also note that it can likewise occur during stages of meditation, deep concentration, sensory deprivation, or with the use of psychedelics such as LSD or mescaline, and even, in some cases, marijuana. However, synesthetes report that common stimulants, like caffeine and cigarettes do not affect the strength of their synesthesia, nor does alcohol.
A very different theoretical approach to synesthesia is that based on ideasthesia. According to this account, synesthesia is a phenomenon mediated by the extraction of the meaning of the inducing stimulus. Thus, synesthesia may be fundamentally a semantic phenomenon. Therefore, to understand neural mechanisms of synesthesia the mechanisms of semantics and the extraction of meaning need to be understood better. This is a non-trivial issue because it is not only a question of a location in the brain at which meaning is "processed" but pertains also to the question of understanding epitomized in e.g., the Chinese room problem. Thus, the question of the neural basis of synesthesia is deeply entrenched into the general mind–body problem and the problem of the explanatory gap.
Genetics
Due to the prevalence of synesthesia among the first-degree relatives of people affected, there may be a genetic basis, as indicated by the monozygotic twins studies showing an epigenetic component. Synesthesia might also be an oligogenic condition, with locus heterogeneity, multiple forms of inheritance, and continuous variation in gene expression. While the exact genetic loci for this trait haven't been identified, research indicates that the genetic constructs underlying synesthesia are most likely more complex than the simple X-linked mode of inheritance that early researchers believed it to be. Further, it remains uncertain as to whether synesthesia perseveres in the genetic pool because it provides a selective advantage, or because it has become a byproduct of some other useful selected trait. Women have a higher chance of developing synesthesia, as demonstrated in population studies conducted in the city of Cambridge, England where females were 6 times more likely to have it. As technological equipment continues to advance, the search for clearer answers regarding the genetics behind synesthesia will become more promising.
Although often termed a "neurological condition," synesthesia is not listed in either the DSM-IV or the ICD since it usually does not interfere with normal daily functioning. Indeed, most synesthetes report that their experiences are neutral or even pleasant. Like perfect pitch, synesthesia is simply a difference in perceptual experience.
The simplest approach is test-retest reliability over long periods of time, using stimuli of color names, color chips, or a computer-screen color picker providing 16.7 million choices. Synesthetes consistently score around 90% on the reliability of associations, even with years between tests. In contrast, non-synesthetes score just 30–40%, even with only a few weeks between tests and a warning that they would be retested.
Many tests exist for synesthesia. Each common type has a specific test. When testing for grapheme–color synesthesia, a visual test is given. The person is shown a picture that includes black letters and numbers. A synesthete will associate the letters and numbers with a specific color. An auditory test is another way to test for synesthesia. A sound is turned on and one will either identify it with a taste or envision shapes. The audio test correlates with chromesthesia (sounds with colors). Since people question whether or not synesthesia is tied to memory, the "retest" is given. One is given a set of objects and is asked to assign colors, tastes, personalities, or more. After some time, the same objects are presented and the person is asked again to do the same task. The synesthete can assign the same characteristics because that person has permanent neural associations in the brain, rather than memories of a certain object.
Grapheme–color synesthetes, as a group, share significant preferences for the color of each letter (e.g., A tends to be red; O tends to be white or black; S tends to be yellow, etc.) Nonetheless, there is a great variety in types of synesthesia, and within each type, individuals report differing triggers for their sensations and differing intensities of experiences. This variety means that defining synesthesia in an individual is difficult, and the majority of synesthetes are completely unaware that their experiences have a name.
Neurologist Richard Cytowic identifies the following diagnostic criteria for synesthesia in his first edition book. However, the criteria are different in the second book:
Synesthesia is involuntary and automatic
Synesthetic perceptions are spatially extended, meaning they often have a sense of "location." For example, synesthetes speak of "looking at" or "going to" a particular place to attend to the experience
Synesthetic percepts are consistent and generic (i.e., simple rather than pictorial)
Synesthesia is highly memorable
Synesthesia is laden with affect
Cytowic's early cases mainly included individuals whose synesthesia was frankly projected outside the body (e.g., on a "screen" in front of one's face). Later research showed that such stark externalization occurs in a minority of synesthetes. Refining this concept, Cytowic and Eagleman differentiated between "localizers" and "non-localizers" to distinguish those synesthetes whose perceptions have a definite sense of spatial quality from those whose perceptions do not.
Prevalence
Estimates of the prevalence of synesthesia have ranged widely, from 1 in 4 to 1 in 25,000–100,000. However, most studies have relied on synesthetes reporting themselves, introducing self-referral bias. In what is cited as the most accurate prevalence study so far, self-referral bias was avoided by studying 500 people recruited from the communities of Edinburgh and Glasgow Universities; it showed a prevalence of 4.4%, with 9 different variations of synesthesia. This study also concluded that one common form of synesthesia grapheme–color synesthesia (colored letters and numbers) is found in more than one percent of the population, and this latter prevalence of graphemes–color synesthesia has since been independently verified in a sample of nearly 3,000 people in the University of Edinburgh.
The most common forms of synesthesia are those that trigger colors, and the most prevalent of all is day–color. Also relatively common is grapheme–color synesthesia. We can think of "prevalence" both in terms of how common is synesthesia (or different forms of synesthesia) within the population, or how common are different forms of synesthesia within synesthetes. So within synesthetes, forms of synesthesia that trigger color also appear to be the most common forms of synesthesia with a prevalence rate of 86% within synesthetes. In another study, music–color is also prevalent at 18–41%. Some of the rarest are reported to be auditory–tactile, mirror-touch, and lexical–gustatory.
There is research to suggest that the likelihood of having synesthesia is greater in people with autism spectrum condition.
History
The interest in colored hearing dates back to Greek antiquity when some theorists wondered whether the color (chroia, what we now call timbre) of music was a quantifiable quality of sound, together with pitch and duration. Additionally, one kind of musical scale (genos) introduced by Plato's friend Archytas of Tarentum in the fourth century BC was named chromatic. The late sixth century BC kitharist Lysander of Sicyon was said to have introduced a more 'colorful' style, even before the development of the chromatic scale itself. In Plato's time, the description of melody as 'colored' had become part of professional jargon, while the musical terms 'tone' and 'harmony' soon became integrated into the vocabulary of color in visual art. Isaac Newton proposed that musical tones and color tones shared common frequencies, as did Goethe in his book Theory of Colours. There is a long history of building color organs such as the clavier à lumières on which to perform colored music in concert halls.
The first medical description of "colored hearing" is in an 1812 thesis by the German physician Georg Tobias Ludwig Sachs. The "father of psychophysics," Gustav Fechner, reported the first empirical survey of colored letter photisms among 73 synesthetes in 1876, followed in the 1880s by Francis Galton. Carl Jung refers to "color hearing" in his Symbols of Transformation in 1912.
In the early 1920s, the Bauhaus teacher and musician Gertrud Grunow researched the relationships between sound, color, and movement and developed a 'twelve-tone circle of colour' which was analogous with the twelve-tone music of the Austrian composer Arnold Schönberg (1874-1951). She was a participant in at least one of the Congresses for Colour-Sound Research (German:Kongreß für Farbe-Ton-Forschung) held in Hamburg in the late 1920s and early 1930s.
Research into synesthesia proceeded briskly in several countries, but due to the difficulties in measuring subjective experiences and the rise of behaviorism, which made the study of any subjective experience taboo, synesthesia faded into scientific oblivion between 1930 and 1980.
As the 1980s cognitive revolution made inquiry into internal subjective states respectable again, scientists returned to studying synesthesia. Led in the United States by Larry Marks and Richard Cytowic, and later in England by Simon Baron-Cohen and Jeffrey Gray, researchers explored the reality, consistency, and frequency of synesthetic experiences. In the late 1990s, the focus settled on grapheme → color synesthesia, one of the most common and easily studied types. Psychologists and neuroscientists study synesthesia not only for its inherent appeal but also for the insights it may give into cognitive and perceptual processes that occur in synesthetes and non-synesthetes alike. Synesthesia is now the topic of scientific books and papers, Ph.D. theses, documentary films, and even novels.
Since the rise of the Internet in the 1990s, synesthetes began contacting one another and creating websites devoted to the condition. These rapidly grew into international organizations such as the American Synesthesia Association, the UK Synaesthesia Association, the Belgian Synesthesia Association, the Canadian Synesthesia Association, the German Synesthesia Association, and the Netherlands Synesthesia Web Community.
Society and culture
Notable cases
Solomon Shereshevsky, a newspaper reporter turned mnemonist, was discovered by Russian neuropsychologist Alexander Luria to have a rare fivefold form of synesthesia, of which he is the only known case. Words and text were not only associated with highly vivid visuospatial imagery but also sound, taste, color, and sensation. Shereshevsky could recount endless details of many things without form, from lists of names to decades-old conversations, but he had great difficulty grasping abstract concepts. The automatic, and nearly permanent, retention of every detail due to synesthesia greatly inhibited Shereshevsky's ability to understand what he read or heard.
Neuroscientist and author V.S. Ramachandran studied the case of a grapheme–color synesthete who was also color blind. While he couldn't see certain colors with his eyes, he could still "see" those colors when looking at certain letters. Because he didn't have a name for those colors, he called them "Martian colors."
Art
Other notable synesthetes come particularly from artistic professions and backgrounds. Synesthetic art historically refers to multi-sensory experiments in the genres of visual music, music visualization, audiovisual art, abstract film, and intermedia. Distinct from neuroscience, the concept of synesthesia in the arts is regarded as the simultaneous perception of multiple stimuli in one gestalt experience. Neurological synesthesia has been a source of inspiration for artists, composers, poets, novelists, and digital artists.
Writers
Vladimir Nabokov wrote explicitly about synesthesia in several novels. Nabokov described his grapheme–color synesthesia at length in his autobiography, Speak, Memory:
I present a fine case of colored hearing. Perhaps "hearing" is not quite accurate, since the color sensations seem to be produced by the very act of my orally forming a given letter while I imagine its outline. The long a of the English alphabet (and it is this alphabet I have in mind farther on unless otherwise stated) has for me the tint of weathered wood, but the French a evokes polished ebony. This black group also includes hard g (vulcanized rubber) and r (a sooty rag being ripped). Oatmeal n, noodle-limp l, and the ivory-backed hand mirror of o take care of the whites. I am puzzled by my French on which I see as the brimming tension-surface of alcohol in a small glass. Passing on to the blue group, there is steely x, thundercloud z, and huckleberry k. Since a subtle interaction exists between sound and shape, I see q as browner than k, while s is not the light blue of c, but a curious mixture of azure and mother-of-pearl.
Daniel Tammet wrote a book on his experiences with synesthesia called Born on a Blue Day. Joanne Harris, author of Chocolat, is a synesthete who says she experiences colors as scents. Her novel Blueeyedboy features various aspects of synesthesia.
Painters and photographers
Wassily Kandinsky (a synesthete) and Piet Mondrian (not a synesthete) both experimented with image–music congruence in their paintings. Contemporary artists with synesthesia, such as Carol Steen and Marcia Smilack (a photographer who waits until she gets a synesthetic response from what she sees and then takes the picture), use their synesthesia to create their artwork. Linda Anderson, according to NPR considered "one of the foremost living memory painters", creates with oil crayons on fine-grain sandpaper representations of the auditory-visual synaesthesia she experiences during severe migraine attacks. Brandy Gale, a Canadian visual artist, experiences an involuntary joining or crossing of any of her senses hearing, vision, taste, touch, smell and movement. Gale paints from life rather than from photographs and by exploring the sensory panorama of each locale attempts to capture, select, and transmit these personal experiences. David Hockney perceives music as color, shape, and configuration and uses these perceptions when painting opera stage sets (though not while creating his other artworks). Kandinsky combined four senses: color, hearing, touch, and smell. American painter and visual artist Perry Hall attributes synaesthesia— both chromesthesia and the experience of visual sensations creating sounds— as an inspiration and guide for his creative work, including his Sound Drawing series.
Composers
Multiple composers had experienced synesthesia.
Mikalojus Konstantinas Čiurlionis, a Lithuanian painter, composer, and writer, perceived colors and music simultaneously. Many of his paintings bear the names of matching musical pieces: sonatas, fugues, and preludes.
Alexander Scriabin composed colored music that was deliberately contrived and based on the circle of fifths, whereas Olivier Messiaen invented a new method of composition (the modes of limited transposition) specifically to render his bi-directional sound–color synesthesia. For example, the red rocks of Bryce Canyon are depicted in his symphony Des canyons aux étoiles... ("From the Canyons to the Stars"). New art movements such as literary symbolism, non-figurative art, and visual music have profited from experiments with synesthetic perception and contributed to the public awareness of synesthetic and multi-sensory ways of perceiving. Other composers who reported synesthesia include Duke Ellington, Nikolay Rimsky-Korsakov, and Jean Sibelius.
Several contemporary composers with a synesthesia are Michael Torke, and Ramin Djawadi, best known for his work on composing the theme songs and scores for such TV series as Game of Thrones, Westworld and for the Iron Man movie. He says he tends to "associate colors with music, or music with colors."
British composer Daniel Liam Glyn created the classical-contemporary music project Changing Stations using Grapheme Colour Synaesthesia. Based on the 11 main lines of the London Underground, the eleven tracks featured on the album represent the eleven main tube line colours. Each track focuses heavily on the different speeds, sounds, and mood of each line, and are composed in the key signature synaesthetically assigned by Glyn with reference to the colour of the tube line on the map.
Musicians
The producer, rapper, and fashion designer Kanye West is a prominent interdisciplinary case. In an impromptu speech he gave during an Ellen interview, he described his condition, saying that he sees sounds and that everything he sonically makes is a painting. Other notable synesthetes include musicians Billy Joel, Andy Partridge, Itzhak Perlman, Lorde, Billie Eilish, Brendon Urie, Ida Maria, Brian Chase, and classical pianist Hélène Grimaud. Musician Kristin Hersh sees music in colors. Drummer Mickey Hart of The Grateful Dead wrote about his experiences with synaesthesia in his autobiography Drumming at the Edge of Magic. John Frusciante, guitarist of the Red Hot Chili Peppers, talks about his experiences with synesthesia in a podcast with Rick Rubin. Pharrell Williams, of the groups The Neptunes and N.E.R.D., also experiences synesthesia and used it as the basis of the album Seeing Sounds. Singer/songwriter Marina and the Diamonds experiences music → color synesthesia and reports colored days of the week. Awsten Knight from Waterparks has chromesthesia, which influences many of the band's artistic choices.
Artists without synesthesia
Some artists frequently mentioned as synesthetes did not, in fact, have the neurological condition. Scriabin's 1911 Prometheus, for example, is a deliberate contrivance whose color choices are based on the circle of fifths and appear to have been taken from Madame Blavatsky. The musical score has a separate staff marked luce whose "notes" are played on a color organ. Technical reviews appear in period volumes of Scientific American. On the other hand, his older colleague Rimsky-Korsakov (who was perceived as a fairly conservative composer) was, in fact, a synesthete.
French poets Arthur Rimbaud and Charles Baudelaire wrote of synesthetic experiences, but there is no evidence they were synesthetes themselves. Baudelaire's 1857 introduced the notion that the senses can and should intermingle. Baudelaire participated in a hashish experiment by psychiatrist Jacques-Joseph Moreau and became interested in how the senses might affect each other. Rimbaud later wrote Voyelles (1871), which was perhaps more important than in popularizing synesthesia. He later boasted "J'inventais la couleur des voyelles!" (I invented the colors of the vowels!).
Science
Some technologists, like inventor Nikola Tesla, and scientists also reported being synesthetic. Physicist Richard Feynman describes his colored equations in his autobiography, What Do You Care What Other People Think?: "When I see equations, I see the letters in colors. I don't know why. I see vague pictures of Bessel functions with light-tan j's, slightly violet-bluish n's, and dark brown x's flying around."
Literature
Synesthesia is sometimes used as a plot device or a way of developing a character's inner life. Author and synesthete Pat Duffy describes four ways in which synesthetic characters have been used in modern fiction.
Synesthesia as Romantic ideal: in which the condition illustrates the Romantic ideal of transcending one's experience of the world. Books in this category include The Gift by Vladimir Nabokov.
Synesthesia as pathology: in which the trait is pathological. Books in this category include The Whole World Over by Julia Glass.
Synesthesia as Romantic pathology: in which synesthesia is pathological but also provides an avenue to the Romantic ideal of transcending quotidian experience. Books in this category include Holly Payne's The Sound of Blue and Anna Ferrara's The Woman Who Tried To Be Normal.
Synesthesia as psychological health and balance: Painting Ruby Tuesday by Jane Yardley, and A Mango-Shaped Space by Wendy Mass.
Literary depictions of synesthesia are criticized as often being more of a reflection of an author's interpretation of synesthesia than of the phenomenon itself.
Research
Research on synesthesia raises questions about how the brain combines information from different sensory modalities, referred to as crossmodal perception or multisensory integration.
An example of this is the bouba/kiki effect. In an experiment first designed by Wolfgang Köhler, people are asked to choose which of two shapes is named bouba and which kiki. The angular shape, kiki, is chosen by 95–98% and bouba for the rounded one. Individuals on the island of Tenerife showed a similar preference between shapes called takete and maluma. Even 2.5-year-old children (too young to read) show this effect. Research indicated that in the background of this effect may operate a form of ideasthesia.
Researchers hope that the study of synesthesia will provide a better understanding of consciousness and its neural correlates. In particular, synesthesia might be relevant to the philosophical problem of qualia, given that synesthetes experience extra qualia (e.g., colored sound). An important insight for qualia research may come from the findings that synesthesia has the properties of ideasthesia, which then suggest a crucial role of conceptualization processes in generating qualia.
Technological applications
Synesthesia also has several practical applications, including 'intentional synesthesia' in technology, and sensory prosthetics.
The Voice (vOICe)
Peter Meijer developed a sensory substitution device for the visually impaired called The vOICe (the capital letters "O," "I," and "C" in "vOICe" are intended to evoke the expression "Oh I see"). The vOICe is a privately owned research project, running without venture capital, that was first implemented using low-cost hardware in 1991. The vOICe is a visual-to-auditory sensory substitution device (SSD) preserving visual detail at high resolution (up to 25,344 pixels). The device consists of a laptop, head-mounted camera or computer camera, and headphones. The vOICe converts visual stimuli of the surroundings captured by the camera into corresponding aural representations (soundscapes) delivered to the user through headphones at a default rate of one soundscape per second. Each soundscape is a left-to-right scan, with height represented by pitch, and brightness by loudness. The vOICe compensates for the loss of vision by converting information from the lost sensory modality into stimuli in a remaining modality.
Sonified
Artists Perry Hall and Jonathan Jones-Morris created Sonified in 2011, software that translates visual information from a video camera into music in real-time, as a means of creating a synaesthetic experience for the user.
| Biology and health sciences | Sensory nervous system | Biology |
21438242 | https://en.wikipedia.org/wiki/Lava | Lava | Lava is molten or partially molten rock (magma) that has been expelled from the interior of a terrestrial planet (such as Earth) or a moon onto its surface. Lava may be erupted at a volcano or through a fracture in the crust, on land or underwater, usually at temperatures from . The volcanic rock resulting from subsequent cooling is often also called lava.
A lava flow is an outpouring of lava during an effusive eruption. (An explosive eruption, by contrast, produces a mixture of volcanic ash and other fragments called tephra, not lava flows.) The viscosity of most lava is about that of ketchup, roughly 10,000 to 100,000 times that of water. Even so, lava can flow great distances before cooling causes it to solidify, because lava exposed to air quickly develops a solid crust that insulates the remaining liquid lava, helping to keep it hot and inviscid enough to continue flowing.
Etymology
The word lava comes from Italian and is probably derived from the Latin word labes, which means a fall or slide. An early use of the word in connection with extrusion of magma from below the surface is found in a short account of the 1737 eruption of Vesuvius, written by Francesco Serao, who described "a flow of fiery lava" as an analogy to the flow of water and mud down the flanks of the volcano (a lahar) after heavy rain.
Properties of lava
Composition
Solidified lava on the Earth's crust is predominantly silicate minerals: mostly feldspars, feldspathoids, olivine, pyroxenes, amphiboles, micas and quartz. Rare nonsilicate lavas can be formed by local melting of nonsilicate mineral deposits or by separation of a magma into immiscible silicate and nonsilicate liquid phases.
Silicate lavas
Silicate lavas are molten mixtures dominated by oxygen and silicon, the most abundant elements of the Earth's crust, with smaller quantities of aluminium, calcium, magnesium, iron, sodium, and potassium and minor amounts of many other elements. Petrologists routinely express the composition of a silicate lava in terms of the weight or molar mass fraction of the oxides of the major elements (other than oxygen) present in the lava.
The silica component dominates the physical behavior of silicate magmas. Silicon ions in lava strongly bind to four oxygen ions in a tetrahedral arrangement. If an oxygen ion is bound to two silicon ions in the melt, it is described as a bridging oxygen, and lava with many clumps or chains of silicon ions connected by bridging oxygen ions is described as partially polymerized. Aluminium in combination with alkali metal oxides (sodium and potassium) also tends to polymerize the lava. Other cations, such as ferrous iron, calcium, and magnesium, bond much more weakly to oxygen and reduce the tendency to polymerize. Partial polymerization makes the lava viscous, so lava high in silica is much more viscous than lava low in silica.
Because of the role of silica in determining viscosity and because many other properties of a lava (such as its temperature) are observed to correlate with silica content, silicate lavas are divided into four chemical types based on silica content: felsic, intermediate, mafic, and ultramafic.
Felsic lava
Felsic or silicic lavas have a silica content greater than 63%. They include rhyolite and dacite lavas. With such a high silica content, these lavas are extremely viscous, ranging from 108 cP (105 Pa⋅s) for hot rhyolite lava at to 1011 cP (108 Pa⋅s) for cool rhyolite lava at . For comparison, water has a viscosity of about 1 cP (0.001 Pa⋅s). Because of this very high viscosity, felsic lavas usually erupt explosively to produce pyroclastic (fragmental) deposits. However, rhyolite lavas occasionally erupt effusively to form lava spines, lava domes or "coulees" (which are thick, short lava flows). The lavas typically fragment as they extrude, producing block lava flows. These often contain obsidian.
Felsic magmas can erupt at temperatures as low as . Unusually hot (>950 °C; >1,740 °F) rhyolite lavas, however, may flow for distances of many tens of kilometres, such as in the Snake River Plain of the northwestern United States.
Intermediate lava
Intermediate or andesitic lavas contain 52% to 63% silica, and are lower in aluminium and usually somewhat richer in magnesium and iron than felsic lavas. Intermediate lavas form andesite domes and block lavas and may occur on steep composite volcanoes, such as in the Andes. They are also commonly hotter than felsic lavas, in the range of . Because of their lower silica content and higher eruptive temperatures, they tend to be much less viscous, with a typical viscosity of 3.5 × 106 cP (3,500 Pa⋅s) at . This is slightly greater than the viscosity of smooth peanut butter. Intermediate lavas show a greater tendency to form phenocrysts. Higher iron and magnesium tends to manifest as a darker groundmass, including amphibole or pyroxene phenocrysts.
Mafic lava
Mafic or basaltic lavas are typified by relatively high magnesium oxide and iron oxide content (whose molecular formulas provide the consonants in mafic) and have a silica content limited to a range of 52% to 45%. They generally erupt at temperatures of and at relatively low viscosities, around 104 to 105 cP (10 to 100 Pa⋅s). This is similar to the viscosity of ketchup, although it is still many orders of magnitude higher than that of water. Mafic lavas tend to produce low-profile shield volcanoes or flood basalts, because the less viscous lava can flow for long distances from the vent. The thickness of a solidified basaltic lava flow, particularly on a low slope, may be much greater than the thickness of the moving molten lava flow at any one time, because basaltic lavas may "inflate" by a continued supply of lava and its pressure on a solidified crust. Most basaltic lavas are of aā or pāhoehoe types, rather than block lavas. Underwater, they can form pillow lavas, which are rather similar to entrail-type pahoehoe lavas on land.
Ultramafic lava
Ultramafic lavas, such as komatiite and highly magnesian magmas that form boninite, take the composition and temperatures of eruptions to the extreme. All have a silica content under 45%. Komatiites contain over 18% magnesium oxide and are thought to have erupted at temperatures of . At this temperature there is practically no polymerization of the mineral compounds, creating a highly mobile liquid. Viscosities of komatiite magmas are thought to have been as low as 100 to 1000 cP (0.1 to 1 Pa⋅s), similar to that of light motor oil. Most ultramafic lavas are no younger than the Proterozoic, with a few ultramafic magmas known from the Phanerozoic in Central America that are attributed to a hot mantle plume. No modern komatiite lavas are known, as the Earth's mantle has cooled too much to produce highly magnesian magmas.
Alkaline lavas
Some silicate lavas have an elevated content of alkali metal oxides (sodium and potassium), particularly in regions of continental rifting, areas overlying deeply subducted plates, or at intraplate hotspots. Their silica content can range from ultramafic (nephelinites, basanites and tephrites) to felsic (trachytes). They are more likely to be generated at greater depths in the mantle than subalkaline magmas. Olivine nephelinite lavas are both ultramafic and highly alkaline, and are thought to have come from much deeper in the mantle of the Earth than other lavas.
Non-silicate lavas
Some lavas of unusual composition have erupted onto the surface of the Earth. These include:
Carbonatite and natrocarbonatite lavas are known from Ol Doinyo Lengai volcano in Tanzania, which is the sole example of an active carbonatite volcano. Carbonatites in the geologic record are typically 75% carbonate minerals, with lesser amounts of silica-undersaturated silicate minerals (such as micas and olivine), apatite, magnetite, and pyrochlore. This may not reflect the original composition of the lava, which may have included sodium carbonate that was subsequently removed by hydrothermal activity, though laboratory experiments show that a calcite-rich magma is possible. Carbonatite lavas show stable isotope ratios indicating they are derived from the highly alkaline silicic lavas with which they are always associated, probably by separation of an immiscible phase. Natrocarbonatite lavas of Ol Doinyo Lengai are composed mostly of sodium carbonate, with about half as much calcium carbonate and half again as much potassium carbonate, and minor amounts of halides, fluorides, and sulphates. The lavas are extremely fluid, with viscosities only slightly greater than water, and are very cool, with measured temperatures of .
Iron oxide lavas are thought to be the source of the iron ore at Kiruna, Sweden which formed during the Proterozoic. Iron oxide lavas of Pliocene age occur at the El Laco volcanic complex on the Chile-Argentina border. Iron oxide lavas are thought to be the result of immiscible separation of iron oxide magma from a parental magma of calc-alkaline or alkaline composition.
Sulfur lava flows up to long and wide occur at Lastarria volcano, Chile. They were formed by the melting of sulfur deposits at temperatures as low as .
The term "lava" can also be used to refer to molten "ice mixtures" in eruptions on the icy satellites of the Solar System's giant planets.
Rheology
The lava's viscosity mostly determines the behavior of lava flows. While the temperature of common silicate lava ranges from about for felsic lavas to for mafic lavas, its viscosity ranges over seven orders of magnitude, from 1011 cP (108 Pa⋅s) for felsic lavas to 104 cP (10 Pa⋅s) for mafic lavas. Lava viscosity is mostly determined by composition but also depends on temperature and shear rate.
Lava viscosity determines the kind of volcanic activity that takes place when the lava is erupted. The greater the viscosity, the greater the tendency for eruptions to be explosive rather than effusive. As a result, most lava flows on Earth, Mars, and Venus are composed of basalt lava. On Earth, 90% of lava flows are mafic or ultramafic, with intermediate lava making up 8% of flows and felsic lava making up just 2% of flows. Viscosity also determines the aspect (thickness relative to lateral extent) of flows, the speed with which flows move, and the surface character of the flows.
When highly viscous lavas erupt effusively rather than in their more common explosive form, they almost always erupt as high-aspect flows or domes. These flows take the form of block lava rather than aā or pāhoehoe. Obsidian flows are common. Intermediate lavas tend to form steep stratovolcanoes, with alternating beds of lava from effusive eruptions and tephra from explosive eruptions. Mafic lavas form relatively thin flows that can move great distances, forming shield volcanoes with gentle slopes.
In addition to melted rock, most lavas contain solid crystals of various minerals, fragments of exotic rocks known as xenoliths, and fragments of previously solidified lava. The crystal content of most lavas gives them thixotropic and shear thinning properties. In other words, most lavas do not behave like Newtonian fluids, in which the rate of flow is proportional to the shear stress. Instead, a typical lava is a Bingham fluid, which shows considerable resistance to flow until a stress threshold, called the yield stress, is crossed. This results in plug flow of partially crystalline lava. A familiar example of plug flow is toothpaste squeezed out of a toothpaste tube. The toothpaste comes out as a semisolid plug, because shear is concentrated in a thin layer in the toothpaste next to the tube and only there does the toothpaste behave as a fluid. Thixotropic behavior also hinders crystals from settling out of the lava. Once the crystal content reaches about 60%, the lava ceases to behave like a fluid and begins to behave like a solid. Such a mixture of crystals with melted rock is sometimes described as crystal mush.
Lava flow speeds vary based primarily on viscosity and slope. In general, lava flows slowly, with typical speeds for Hawaiian basaltic flows of and maximum speeds of on steep slopes. An exceptional speed of was recorded following the collapse of a lava lake at Mount Nyiragongo. The scaling relationship for lavas is that the average speed of a flow scales as the square of its thickness divided by its viscosity. This implies that a rhyolite flow would have to be about a thousand times thicker than a basalt flow to flow at a similar speed.
Temperature
The temperature of most types of molten lava ranges from about to depending on the lava's chemical composition. This temperature range is similar to the hottest temperatures achievable with a forced air charcoal forge. Lava is most fluid when first erupted, becoming much more viscous as its temperature drops.
Lava flows quickly develop an insulating crust of solid rock as a result of radiative loss of heat. Thereafter, the lava cools by a very slow conduction of heat through the rocky crust. For instance, geologists of the United States Geological Survey regularly drilled into the Kilauea Iki lava lake, formed in an eruption in 1959. After three years, the solid surface crust, whose base was at a temperature of , was still only thick, even though the lake was about deep. Residual liquid was still present at depths of around nineteen years after the eruption.
A cooling lava flow shrinks, and this fractures the flow. Basalt flows show a characteristic pattern of fractures. The uppermost parts of the flow show irregular downward-splaying fractures, while the lower part of the flow shows a very regular pattern of fractures that break the flow into five- or six-sided columns. The irregular upper part of the solidified flow is called the entablature, while the lower part that shows columnar jointing is called the colonnade. (The terms are borrowed from Greek temple architecture.) Likewise, regular vertical patterns on the sides of columns, produced by cooling with periodic fracturing, are described as chisel marks. Despite their names, these are natural features produced by cooling, thermal contraction, and fracturing.
As lava cools, crystallizing inwards from its edges, it expels gases to form vesicles at the lower and upper boundaries. These are described as pipe-stem vesicles or pipe-stem amygdales. Liquids expelled from the cooling crystal mush rise upwards into the still-fluid center of the cooling flow and produce vertical vesicle cylinders. Where these merge towards the top of the flow, they form sheets of vesicular basalt and are sometimes capped with gas cavities that sometimes fill with secondary minerals. The beautiful amethyst geodes found in the flood basalts of South America formed in this manner.
Flood basalts typically crystallize little before they cease flowing, and, as a result, flow textures are uncommon in less silicic flows. On the other hand, flow banding is common in felsic flows.
Lava morphology
The morphology of lava describes its surface form or texture. More fluid basaltic lava flows tend to form flat sheet-like bodies, whereas viscous rhyolite lava flows form knobbly, blocky masses of rock. Lava erupted underwater has its own distinctive characteristics.
Aā
Aā (also spelled aa, aa, aa, and a-aa, and pronounced or ) is one of three basic types of flow lava. Aā is basaltic lava characterized by a rough or rubbly surface composed of broken lava blocks called clinker. The word is Hawaiian meaning "stony rough lava", but also to "burn" or "blaze"; it was introduced as a technical term in geology by Clarence Dutton.
The loose, broken, and sharp, spiny surface of an aā flow makes hiking difficult and slow. The clinkery surface actually covers a massive dense core, which is the most active part of the flow. As pasty lava in the core travels downslope, the clinkers are carried along at the surface. At the leading edge of an aā flow, however, these cooled fragments tumble down the steep front and are buried by the advancing flow. This produces a layer of lava fragments both at the bottom and top of an aā flow.
Accretionary lava balls as large as are common on aā flows. Aā is usually of higher viscosity than pāhoehoe. Pāhoehoe can turn into aā if it becomes turbulent from meeting impediments or steep slopes.
The sharp, angled texture makes aā a strong radar reflector, and can easily be seen from an orbiting satellite (bright on Magellan pictures).
Aā lavas typically erupt at temperatures of or greater.
Pāhoehoe
Pāhoehoe (also spelled pahoehoe, from Hawaiian meaning "smooth, unbroken lava") is basaltic lava that has a smooth, billowy, undulating, or ropy surface. These surface features are due to the movement of very fluid lava under a congealing surface crust. The Hawaiian word was introduced as a technical term in geology by Clarence Dutton.
A pāhoehoe flow typically advances as a series of small lobes and toes that continually break out from a cooled crust. It also forms lava tubes where the minimal heat loss maintains a low viscosity. The surface texture of pāhoehoe flows varies widely, displaying all kinds of bizarre shapes often referred to as lava sculpture. With increasing distance from the source, pāhoehoe flows may change into aā flows in response to heat loss and consequent increase in viscosity. Experiments suggest that the transition takes place at a temperature between , with some dependence on shear rate. Pahoehoe lavas typically have a temperature of .
On the Earth, most lava flows are less than long, but some pāhoehoe flows are more than long. Some flood basalt flows in the geologic record extend for hundreds of kilometres.
The rounded texture makes pāhoehoe a poor radar reflector, and is difficult to see from an orbiting satellite (dark on Magellan picture).
Block lava flows
Block lava flows are typical of andesitic lavas from stratovolcanoes. They behave in a similar manner to aā flows but their more viscous nature causes the surface to be covered in smooth-sided angular fragments (blocks) of solidified lava instead of clinkers. As with aā flows, the molten interior of the flow, which is kept insulated by the solidified blocky surface, advances over the rubble that falls off the flow front. They also move much more slowly downhill and are thicker in depth than aā flows.
Pillow lava
Pillow lava is the lava structure typically formed when lava emerges from an underwater volcanic vent or subglacial volcano or a lava flow enters the ocean. The viscous lava gains a solid crust on contact with the water, and this crust cracks and oozes additional large blobs or "pillows" as more lava emerges from the advancing flow. Since water covers the majority of Earth's surface and most volcanoes are situated near or under bodies of water, pillow lava is very common.
Lava landforms
Because it is formed from viscous molten rock, lava flows and eruptions create distinctive formations, landforms and topographical features from the macroscopic to the microscopic.
Volcanoes
Volcanoes are the primary landforms built by repeated eruptions of lava and ash over time. They range in shape from shield volcanoes with broad, shallow slopes formed from predominantly effusive eruptions of relatively fluid basaltic lava flows, to steeply-sided stratovolcanoes (also known as composite volcanoes) made of alternating layers of ash and more viscous lava flows typical of intermediate and felsic lavas.
A caldera, which is a large subsidence crater, can form in a stratovolcano, if the magma chamber is partially or wholly emptied by large explosive eruptions; the summit cone no longer supports itself and thus collapses in on itself afterwards. Such features may include volcanic crater lakes and lava domes after the event. However, calderas can also form by non-explosive means such as gradual magma subsidence. This is typical of many shield volcanoes.
Cinder and spatter cones
Cinder cones and spatter cones are small-scale features formed by lava accumulation around a small vent on a volcanic edifice. Cinder cones are formed from tephra or ash and tuff which is thrown from an explosive vent. Spatter cones are formed by accumulation of molten volcanic slag and cinders ejected in a more liquid form.
Kīpukas
Another Hawaiian English term derived from the Hawaiian language, a kīpuka denotes an elevated area such as a hill, ridge or old lava dome inside or downslope from an area of active volcanism. New lava flows will cover the surrounding land, isolating the kīpuka so that it appears as a (usually) forested island in a barren lava flow.
Lava domes and coulées
Lava domes are formed by the extrusion of viscous felsic magma. They can form prominent rounded protuberances, such as at Valles Caldera. As a volcano extrudes silicic lava, it can form an inflation dome or endogenous dome, gradually building up a large, pillow-like structure which cracks, fissures, and may release cooled chunks of rock and rubble. The top and side margins of an inflating lava dome tend to be covered in fragments of rock, breccia and ash.
Examples of lava dome eruptions include the Novarupta dome, and successive lava domes of Mount St Helens.
When a dome forms on an inclined surface it can flow in short thick flows called coulées (dome flows). These flows often travel only a few kilometres from the vent.
Lava tubes
Lava tubes are formed when a flow of relatively fluid lava cools on the upper surface sufficiently to form a crust. Beneath this crust, which being made of rock is an excellent insulator, the lava can continue to flow as a liquid. When this flow occurs over a prolonged period of time the lava conduit can form a tunnel-like aperture or lava tube, which can conduct molten rock many kilometres from the vent without cooling appreciably. Often these lava tubes drain out once the supply of fresh lava has stopped, leaving a considerable length of open tunnel within the lava flow.
Lava tubes are known from the modern day eruptions of Kīlauea, and significant, extensive and open lava tubes of Tertiary age are known from North Queensland, Australia, some extending for .
Lava lakes
Rarely, a volcanic cone may fill with lava but not erupt. Lava which pools within the caldera is known as a lava lake. Lava lakes do not usually persist for long, either draining back into the magma chamber once pressure is relieved (usually by venting of gases through the caldera), or by draining via eruption of lava flows or pyroclastic explosion.
There are only a few sites in the world where permanent lakes of lava exist. These include:
Mount Erebus, Antarctica
Erta Ale, Ethiopia
Nyiragongo, Democratic Republic of Congo
Ambrym, Vanuatu.
Lava delta
Lava deltas form wherever sub-aerial flows of lava enter standing bodies of water. The lava cools and breaks up as it encounters the water, with the resulting fragments filling in the seabed topography such that the sub-aerial flow can move further offshore. Lava deltas are generally associated with large-scale, effusive type basaltic volcanism.
Lava fountains
A lava fountain is a volcanic phenomenon in which lava is forcefully but non-explosively ejected from a crater, vent, or fissure. The highest lava fountain recorded was during the 23 November 2013 eruption of Mount Etna in Italy, which reached a stable height of around for 18 minutes, briefly peaking at a height of . Lava fountains may occur as a series of short pulses, or a continuous jet of lava. They are commonly associated with Hawaiian eruptions.
Hazards
Lava flows are enormously destructive to property in their path. However, casualties are rare since flows are usually slow enough for people and animals to escape, though this is dependent on the viscosity of the lava. Nevertheless, injuries and deaths have occurred, either because they had their escape route cut off, because they got too close to the flow or, more rarely, if the lava flow front travels too quickly. This notably happened during the eruption of Nyiragongo in Zaire (now Democratic Republic of the Congo). On the night of 10 January 1977, a crater wall was breached and a fluid lava lake drained out in under an hour. The resulting flow sped down the steep slopes at up to , and overwhelmed several villages while residents were asleep. As a result of this disaster, the mountain was designated a Decade Volcano in 1991.
Deaths attributed to volcanoes frequently have a different cause. For example, volcanic ejecta, pyroclastic flow from a collapsing lava dome, lahars, poisonous gases that travel ahead of lava, or explosions caused when the flow comes into contact with water. A particularly dangerous area is called a lava bench. This very young ground will typically break off and fall into the sea.
Areas of recent lava flows continue to represent a hazard long after the lava has cooled. Where young flows have created new lands, land is more unstable and can break off into the sea. Flows often crack deeply, forming dangerous chasms, and a fall against aā lava is similar to falling against broken glass. Rugged hiking boots, long pants, and gloves are recommended when crossing lava flows.
Diverting a lava flow is extremely difficult, but it can be accomplished in some circumstances, as was once partially achieved in Vestmannaeyjar, Iceland. The optimal design of simple, low-cost barriers that divert lava flows is an area of ongoing research.
Towns destroyed by lava flows
The Nisga'a villages of Lax Ksiluux and Wii Lax K'abit in northwestern British Columbia, Canada, were destroyed by thick lava flows during the eruption of Tseax Cone in the 1700s.
Garachico on the island of Tenerife was destroyed by the eruption of Trevejo (1706) (rebuilt).
Cagsawa, Philippines, buried by lava erupted from Mayon Volcano in 1814
Keawaiki, Hawaii, 1859 (abandoned)
San Sebastiano al Vesuvio, Italy, destroyed in 1944 by the most recent eruption of Mount Vesuvius during the Allies' occupation of southern Italy (rebuilt)
Koae and Kapoho, Hawaii, were both destroyed by the same eruption of Kīlauea in January, 1960 (abandoned).
Kalapana, Hawaii, was destroyed by the eruption of the Kīlauea volcano in 1990 (abandoned).
Kapoho, Hawaii, was largely inundated by lava in June 2018, with its subdivision Vacationland Hawaii being completely destroyed.
Towns damaged by lava flows
Catania, Italy, in the 1669 Etna eruption (rebuilt)
Sale'aula, Samoa, by eruptions of Mt Matavanu between 1905 and 1911
Mascali, Italy, almost destroyed by the eruption of Mount Etna in 1928 (rebuilt)
Parícutin (village after which the volcano was named) and San Juan Parangaricutiro, Mexico, by Parícutin from 1943 to 1952
Heimaey, Iceland, in the 1973 Eldfell eruption (rebuilt)
Piton Sainte-Rose, Reunion island, in 1977
Royal Gardens, Hawaii, by the eruption of Kilauea in 1986–87 (abandoned)
Goma, Democratic Republic of Congo, in the eruption of Nyiragongo in 2002
Los Llanos de Aridane (Todoque neighbourhood) and El Paso (El Paraíso neighbourhood) on La Palma in the 2021 Cumbre Vieja volcanic eruption
Towns destroyed by tephra
Tephra is lava in the form of volcanic ash, lapilli, volcanic bombs or volcanic blocks.
Pompeii, Italy, in the eruption of Mount Vesuvius in 79 AD
Herculaneum, Italy, in the eruption of Mount Vesuvius in 79 AD
Cerén, El Salvador, in the eruption of Ilopango between 410 and 535 AD
Sumbawa Island, Indonesia, in the eruption of Mount Tambora in 1815
Plymouth, Montserrat, in 1995. Plymouth was the capital and only port of entry for Montserrat and had to be completely abandoned, along with over half of the island. It is still the de jure capital.
| Physical sciences | Volcanic landforms | null |
21438807 | https://en.wikipedia.org/wiki/Oncology | Oncology | Oncology is a branch of medicine that deals with the study, treatment, diagnosis, and prevention of cancer. A medical professional who practices oncology is an oncologist. The name's etymological origin is the Greek word ὄγκος (ónkos), meaning "tumor", "volume" or "mass". Oncology is concerned with:
The diagnosis of any cancer in a person (pathology)
Therapy (e.g. surgery, chemotherapy, radiotherapy and other modalities)
Follow-up of cancer patients after successful treatment
Palliative care of patients with terminal malignancies
Ethical questions surrounding cancer care
Screening efforts:
of populations, or
of the relatives of patients (in types of cancer that are thought to have a hereditary basis, such as breast cancer)
Diagnosis
Medical histories remain an important screening tool: the character of the complaints and nonspecific symptoms (such as fatigue, weight loss, unexplained anemia, fever of unknown origin, paraneoplastic phenomena and other signs) may warrant further investigation for malignancy. Occasionally, a physical examination may find the location of a malignancy.
Diagnostic methods include:
Biopsy or resection; these are methods by which suspicious neoplastic growths can be removed in part or in whole, and evaluated by a pathologist to determine malignancy. This is currently the gold standard for the diagnosis of cancer and is crucial in guiding the next step in management (active surveillance, surgery, radiation therapy, chemotherapy, or a combination of these)
Endoscopy, either upper or lower gastrointestinal, cystoscopy, bronchoscopy, or nasendoscopy; to localise areas suspicious for malignancy and biopsy when necessary.
Mammograms, X-rays, CT scanning, MRI scanning, ultrasound and other radiological techniques to localise and guide biopsy.
Scintigraphy, single photon emission computed tomography (SPECT), positron emission tomography (PET) and other methods of nuclear medicine to identify areas suspicious of malignancy.
Blood tests, including tumor markers, which can increase the suspicion of certain types of cancers.
Apart from diagnoses, these modalities (especially imaging by CT scanning) are often used to determine operability, i.e. whether it is surgically possible to remove a tumor in its entirety.
Currently, a tissue diagnosis (from a biopsy) by a pathologist is essential for the proper classification of cancer and to guide the next step of treatment. On extremely rare instances when this is not possible, "empirical therapy" (without an exact diagnosis) may be considered, based on the available evidence (e.g. history, x-rays and scans.)
On very rare occasions, a metastatic lump or pathological lymph node is found (typically in the neck) for which a primary tumor cannot be found. However, immunohistochemical markers often give a strong indication of the primary malignancy. This situation is referred to as "malignancy of unknown primary", and again, treatment is empirically based on past experience of the most likely origin.
Therapy
Depending upon the cancer identified, follow-up and palliative care will be administered at that time. Certain disorders (such as ALL or AML) will require immediate admission and chemotherapy, while others will be followed up with regular physical examination and blood tests.
Often, surgery is attempted to remove a tumor entirely. This is only feasible when there is some degree of certainty that the tumor can in fact be removed. When it is certain that parts will remain, curative surgery is often impossible, e.g. when there are metastases, or when the tumor has invaded a structure that cannot be operated upon without risking the patient's life. Occasionally surgery can improve survival even if not all tumour tissue has been removed; the procedure is referred to as "debulking" (i.e. reducing the overall amount of tumour tissue). Surgery is also used for the palliative treatment of some cancers, e.g. to relieve biliary obstruction, or to relieve the problems associated with some cerebral tumors. The risks of surgery must be weighed against the benefits.
Chemotherapy and radiotherapy are used as a first-line radical therapy in several malignancies. They are also used for adjuvant therapy, i.e. when the macroscopic tumor has already been completely removed surgically but there is a reasonable statistical risk that it will recur. Chemotherapy and radiotherapy are commonly used for palliation, where disease is clearly incurable: in this situation the aim is to improve the quality of life and to prolong it.
Hormone manipulation is well established, particularly in the treatment of breast and prostate cancer.
There is currently a rapid expansion in the use of monoclonal antibody treatments, notably for lymphoma (Rituximab) and breast cancer (Trastuzumab).
Vaccines and other immunotherapies are the subject of intensive research.
Palliative care
Approximately 50% of all cancer cases in the Western world can be treated to remission with radical treatment. For pediatric patients, that number is much higher. A large number of cancer patients will die from the disease, and a significant proportion of patients with incurable cancer will die of other causes. There may be ongoing issues with symptom control associated with progressive cancer, and also with the treatment of the disease. These problems may include pain, nausea, anorexia, fatigue, immobility, and depression. Not all issues are strictly physical: personal dignity may be affected. Moral and spiritual issues are also important.
While many of these problems fall within the remit of the oncologist, palliative care has matured into a separate, closely allied specialty to address the problems associated with advanced disease. Palliative care is an essential part of the multidisciplinary cancer care team. Palliative care services may be less hospital-based than oncology, with nurses and doctors who are able to visit the patient at home.
Ethical issues
There are a number of recurring ethical questions and dilemmas in oncological practice. These include:
What information to give the patient regarding disease extent/progression/prognosis.
Entry into clinical trials, especially in the face of terminal illness.
Withdrawal of active treatment.
"Do Not Resuscitate" orders and other end-of-life issues.
These issues are closely related to the patient's personality, religion, culture, and family life. Though these issues are complex and emotional, the answers are often achieved by the patient seeking counsel from trusted personal friends and advisors. It requires a degree of sensitivity and very good communication on the part of the oncology team to address these problems properly.
Progress and research
There is a tremendous amount of research being conducted on all frontiers of oncology, ranging from cancer cell biology, and radiation therapy to chemotherapy treatment regimens and optimal palliative care and pain relief. Next-generation sequencing and whole-genome sequencing have completely changed the understanding of cancers. Identification of novel genetic/molecular markers will change the methods of diagnosis and treatment, paving the way for personalized medicine.
Therapeutic trials often involve patients from many different hospitals in a particular region. In the UK, patients are often enrolled in large studies coordinated by Cancer Research UK (CRUK), Medical Research Council (MRC), the European Organisation for Research and Treatment of Cancer (EORTC) or the National Cancer Research Network (NCRN).
The most valued companies worldwide whose leading products are in Oncology include Pfizer (United States), Roche (Switzerland), Merck (United States), AstraZeneca (United Kingdom), Novartis (Switzerland) and Bristol-Myers Squibb (United States) who are active in the treatment areas Kinase inhibitors, Antibodies, Immuno-oncology and Radiopharmaceuticals.
Specialties
The four main divisions:
Medical oncology: focuses on the treatment of cancer with chemotherapy, targeted therapy, immunotherapy, and hormonal therapy.
Surgical oncology: focuses on treatment of cancer with surgery.
Radiation oncology: focuses on treatment of cancer with radiation.
Clinical oncology: focuses on treatment of cancer with both systemic therapies and radiation.
Sub-specialties in Oncology:
Neuro-oncology: focuses on cancers of brain.
Ocular oncology: focuses on cancers of eye.
Head & Neck oncology: focuses on cancers of oral cavity, nasal cavity, oropharynx, hypopharynx and larynx.
Thoracic oncology: focuses on cancers of lung, mediastinum, oesophagus and pleura.
Breast oncology: focuses on cancers of breast.
Gastrointestinal oncology: focuses on cancers of the stomach, colon, rectum, anal canal, liver, gallbladder, pancreas.
Bone & Musculoskeletal oncology: focuses on cancers of bones and soft tissue.
Dermatological oncology: focuses on the medical and surgical treatment of skin, hair, sweat gland, and nail cancers
Genitourinary oncology: focuses on cancers of genital and urinary system.
Gynecologic oncology: focuses on cancers of the female reproductive system.
Pediatric oncology: concerned with the treatment of cancer in children.
Adolescent and young adult (AYA) oncology.
Hemato oncology: focuses on cancers of blood and stem cell transplantation.
Preventive oncology: focuses on epidemiology & prevention of cancer.
Geriatric oncology: focuses on cancers in elderly population.
Pain & Palliative oncology: focuses on treatment of end stage cancer to help alleviate pain and suffering.
Molecular oncology: focuses on molecular diagnostic methods in oncology.
Nuclear medicine oncology: focuses on diagnosis and treatment of cancer with radiopharmaceuticals.
Psycho-oncology: focuses on psychosocial issues on diagnosis and treatment of cancer patients.
Veterinary oncology: focuses on treatment of cancer in animals.
Emerging specialties:
Cardiooncology is a branch of cardiology that addresses the cardiovascular impact of cancer and its treatments.
| Biology and health sciences | Fields of medicine | null |
2858897 | https://en.wikipedia.org/wiki/Demersal%20fish | Demersal fish | Demersal fish, also known as groundfish, live and feed on or near the bottom of seas or lakes (the demersal zone). They occupy the sea floors and lake beds, which usually consist of mud, sand, gravel or rocks. In coastal waters, they are found on or near the continental shelf, and in deep waters, they are found on or near the continental slope or along the continental rise. They are not generally found in the deepest waters, such as abyssal depths or on the abyssal plain, but they can be found around seamounts and islands. The word demersal comes from the Latin demergere, which means to sink.
Demersal fish are bottom feeders. They can be contrasted with pelagic fish, which live and feed away from the bottom in the open water column.
Demersal fish fillets contain little fish oil (one to four per cent), whereas pelagic fish can contain up to 30 per cent.
Types
Demersal fish can be divided into two main types: strictly benthic fish which can rest on the sea floor, and benthopelagic fish which can float in the water column just above the sea floor.
Benthopelagic fish have neutral buoyancy, so they can float at depth without much effort, while strictly benthic fish are denser, with negative buoyancy so they can lie on the bottom without any effort. Most demersal fish are benthopelagic.
As with other bottom feeders, a mechanism to deal with substrate is often necessary. With demersal fish the sand is usually pumped out of the mouth through the gill slit. Most demersal fish exhibit a flat ventral region so as to more easily rest their body on the substrate. The exception may be the flatfish, which are laterally depressed but lie on their sides. Also, many exhibit what is termed an "inferior" mouth, which means that the mouth is pointed downwards; this is beneficial as their food is often below them in the substrate. Those bottom feeders with upward-pointing mouths, such as stargazers, tend to seize swimming prey.
Benthic fish
Benthic fish are denser than water, so they can rest on the sea floor. They either lie-and-wait as ambush predators, at times covering themselves with sand or otherwise camouflaging themselves, or move actively over the bottom in search for food. Benthic fish which can bury themselves include dragonets, flatfish and stingrays.
Flatfish are an order of ray-finned benthic fishes which lie flat on the ocean floor. Examples are flounder, sole, turbot, plaice, and halibut. The adult fish of many species have both eyes on one side of the head. When the fish hatches, one eye is located on each side of its head. But as the fish grows from the larval stage, one eye migrates to the other side of the body as a process of metamorphosis. The flatfish then changes its habits, and camouflages itself by lying on the bottom of the ocean floor with both eyes facing upwards. The side to which one eye migrates depends on the species; with some species both eyes are ultimately on the left side, whereas with other species the eyes are on the right.
Flounder ambush their prey, feeding at soft muddy area of the sea bottom, near bridge piles, docks, artificial and coral reefs. Their diet consists mainly of fish spawn, crustaceans, polychaetes and small fish.
The great hammerhead swings its head in broad angles over the sea floor to pick up the electrical signatures of stingrays buried in the sand. It then uses its "hammer" to pin down the stingray.
Some fishes do not fit into the above classification. For example, the family of nearly blind spiderfishes, common and widely distributed, feed on benthopelagic zooplankton. Yet they are strictly benthic fish, since they stay in contact with the bottom. Their fins have long rays they use to "stand" on the bottom while they face the current and grab zooplankton as it passes by.
The bodies of benthic fish are adapted for ongoing contact with the sea floor. Swimbladders are usually absent or reduced, and the bodies are usually flattened in one way or another. Following Moyle and Cech (2004) they can be divided into five overlapping body shapes:
Benthopelagic fish
Benthopelagic fish inhabit the water just above the bottom, feeding on benthos and zooplankton. Most demersal fish are benthopelagic.
Deep sea benthopelagic teleosts all have swimbladders. The dominant species, rattails and cusk eels, have considerable biomass. Other species include deep sea cods (morids), deep sea eels, halosaurs and notacanths.
Benthopelagic sharks, like the deep sea squaloid sharks, achieve neutral buoyancy with the use of large oil-filled livers. Sharks adapt well to fairly high pressures. They can often be found on slopes down to about 2000 metres, scavenging on food falls such as dead whales. However, the energy demands of sharks are high, since they need to swim constantly and maintain a large amount of oil for buoyancy. These energy needs cannot be met in the extreme oligotrophic conditions that occur at great depths.
Shallow water stingrays are benthic, and can lie on the bottom because of their negative buoyancy. Deep sea stingrays are benthopelagic, and like the squaloids have very large livers which give them neutral buoyancy.
Benthopelagic fish can be divided into flabby or robust body types. Flabby benthopelagic fishes are like bathypelagic fishes; they have a reduced body mass, and low metabolic rates, expending minimal energy as they lie and wait to ambush prey. An example of a flabby fish is the cusk-eel Acanthonus armatus, a predator with a huge head and a body that is 90 per cent water. This fish has the largest ears (otoliths) and the smallest brain in relation to its body size of all known vertebrates.
Deepwater benthopelagic fish are robust, muscular swimmers that actively cruise the bottom searching for prey. They often live around features, such as seamounts, which have strong currents. Commercial examples are the orange roughy and Patagonian toothfish.
Habitats
The edge of the continental shelf marks the boundary where the shelf gives way to, and then gradually drops into abyssal depths. This edge marks the boundary between coastal, relatively shallow, benthic habitats, and the deep benthic habitats. Coastal demersal fishes live on the bottom of inshore waters, such as bays and estuaries, and further out, on the floor of the continental shelf. Deep water demersal fish live beyond this edge, mostly down the continental slopes and along the continental rises which drop to the abyssal plains. This is the continental margin, constituting about 28% of the total oceanic area. Other deep sea demersal fish can also be found around seamounts and islands.
The term bathydemersal fish is sometimes used instead of "deep water demersal fish". Bathydemersal refers to demersal fish which live at depths greater than 200 metres.
The term epibenthic is also used to refer to organism that live on top of the ocean floor, as opposed to those that burrow into the seafloor substrate. However the terms mesodemersal, epidemersal, mesobenthic and bathybenthic are not used.
Coastal
Coastal demersal fish are found on or near the seabed of coastal waters between the shoreline and the edge of the continental shelf, where the shelf drops into the deep ocean. Since the continental shelf is generally less than 200 metres deep, this means that coastal waters are generally epipelagic. The term includes demersal reef fish and demersal fish that inhabit estuaries, inlets and bays.
Young mangrove jacks, a sought after eating and sport fish, dwell in estuaries around mangrove roots, fallen trees, rock walls, and any other snag areas where smaller prey reside for protection. When they mature, they migrate into open waters, sometimes hundreds of kilometres from the coast to spawn.
Stargazers are found worldwide in shallow waters. They have eyes on top of their heads and a large upward-facing mouth. They bury themselves in sand, and leap upwards to ambush benthopelagic fish and invertebrates that pass overhead. Some species have a worm-shaped lure growing out of the floor of the mouth, which they wiggle to attract prey. Stargazers are venomous and can deliver electric shocks. They have been called "the meanest things in creation."
Other examples of coastal demersal fish are cod, plaice, monkfish and sole.
Deep water
Deep water demersal fish occupy the benthic regions beyond the continental margins.
On the continental slope, demersal fishes are common. They are more diverse than coastal demersal fish, since there is more habitat diversity. Further out are the abyssal plains. These flat, featureless regions occupy about 40 per cent of the ocean floor. They are covered with sediment but largely devoid of benthic life (benthos). Deep sea benthic fishes are more likely to associate with canyons or rock outcroppings among the plains, where invertebrate communities are established. Undersea mountains (seamounts) can intercept deep sea currents, and cause productive upwellings which support benthic fish. Undersea mountain ranges can separate underwater regions into different ecosystems.
Rattails and brotulas are common, and other well-established families are eels, eelpouts, hagfishes, greeneyes, batfishes and lumpfishes.
The bodies of deep water demersal fishes are muscular with well developed organs. In this way they are closer to mesopelagic fishes than bathypelagic fishes. In other ways, they are more variable. Photophores are usually absent, eyes and swimbladders range from absent to well developed. They vary in size, and larger species, greater than one metre, are not uncommon.
Deep sea demersal fish are usually long and narrow. Many are eels or shaped like eels. This may be because long bodies have long lateral lines. Lateral lines detect low-frequency sounds, and some demersal fishes have muscles that drum such sounds to attract mates. Smell is also important, as indicated by the rapidity with which demersal fish find traps baited with bait fish.
The main diet of deep sea demersal fish is invertebrates of the deep sea benthos and carrion. Smell, touch and lateral line sensitivities seem to be the main sensory devices for locating these.
Like coastal demersal fish, deep sea demersal fish can be divided into benthic fish and benthopelagic fish, where the benthic fish are negatively buoyant and benthopelagic fish are neutrally buoyant.
The availability of plankton for food diminishes rapidly with depth. At , the biomass of plankton is typically about 1 per cent of that at the surface, and at about 0.01 per cent. Given there is no sunlight, energy enters deep water zones as organic matter. There are three main ways this happens. Firstly, organic matter can move into the zone from the continental landmass, for example, through currents that carry the matter down rivers, then plume along the continental shelf and finally spill down the continental slope. Other matter enters as particulate matter raining down from the overhead water column in the form of marine snow, or as sinking overhead plant material such as eelgrass, or as "large particles" such as dead fish and whales sinking to the bottom. A third way energy can arrive is through fish, such as vertically migrating mesopelagic fishes that can enter into the demersal zone as they ascend or descend. The demersal fish and invertebrates consume organic matter that does arrive, break it down and recycle it. A consequence of these energy delivery mechanisms is that the abundance of demersal fish and invertebrates gradually decrease as the distance from continental shorelines increases.
Although deep water demersal fish species are not generally picky about what they eat, there is still some degree of specialisation. For example, different fish have different mouth sizes, which determines the size of the prey they can handle. Some feed mostly on benthopelagic organisms. Others fed mostly on epifauna (invertebrates on top of the seafloor surface, also called epibenthos), or alternatively on infauna (invertebrates that burrow into the seafloor substrate). Infauna feeders can have considerable sediment in their stomachs. Scavengers, such as snubnosed eels and hagfish, also eat infauna as a secondary food source.
Some feed on carrion. Cameras show that when a dead fish is placed on the bottom, vertebrate and invertebrate scavengers appear very quickly. If the fish is large, some scavengers burrow in and eat it from the inside out. Some fish, such as grenadiers, also appear and start feeding on the scavenging invertebrates and amphipods. Other specialization is based on depth distribution. Some of the more abundant upper continental slope fish species, such as cutthroat eel and longfinned hake, mainly feed on epipelagic fish. But generally, the most abundant deep water demersal fish species feed on invertebrates.
At great depths, food scarcity and extreme pressure limits the ability of fish to survive. The deepest point of the ocean is about 11,000 metres. Bathypelagic fishes are not normally found below 3,000 metres. It may be that extreme pressures interfere with essential enzyme functions.
The deepest-living fish known, the strictly benthic Abyssobrotula galatheae, eel-like and blind, feeds on benthic invertebrates. A living example was trawled from the bottom of the Puerto Rico Trench in 1970 from a depth of 8,370 metres (27,453 ft).
In 2008, a shoal of 17 hadal snailfish, a species of deep water snailfish, was filmed by a UK-Japan team using remote operated landers at depths of in the Japan Trench in the Pacific. The fish were 30 centimetres long (12 in), and were darting about, using vibration sensors on their nose to catch shrimps. The team also reported that the appearance of the fish, unlike that of most deep sea fish, was surprisingly "cute", and that they were surprised by how active the fish were at these depths.
Demersal fisheries
Most demersal fish of commercial or recreational interest are coastal, confined to the upper 200 metres. Commercially important demersal food fish species include flatfish, such as flounder, sole, turbot, plaice, and halibut. Also important are cod, hake, redfish, haddock, bass, congers, sharks, rays and chimaeras.
The following table shows the world capture production of some groups of demersal species in tonnes.
Black sea bass inhabit US coasts from Maine to NE Florida and the eastern Gulf of Mexico, and are most abundant off the waters of New York. They are found in inshore waters (bays and sounds) and offshore in waters up to a depth of 130 m (425'). They spend most of their time close to the sea floor and are often congregated around bottom formations such as rocks, man-made reefs, wrecks, jetties, piers, and bridge pilings. Black sea bass are sought after recreational and commercial fish, and have been overfished.
Grouper are often found around reefs. They have stout bodies and large mouths. They are not built for long-distance or fast swimming. They can be quite large, and lengths over a meter and weights up to 100 kg are not uncommon. They swallow prey rather than biting pieces off it. They do not have many teeth on the edges of their jaws, but they have heavy crushing tooth plates inside the pharynx. They lie in wait, rather than chasing in open water. They are found in areas of hard or consolidated substrate, and use structural features such as ledges, rocks, and coral reefs (as well as artificial reefs like wrecks and sunken barges) as their habitat. Their mouth and gills form a powerful sucking system that sucks their prey in from a distance. They also use their mouth to dig into sand to form their shelters under big rocks, jetting it out through their gills. Their gill muscles are so powerful that it is nearly impossible to pull them out of their cave if they feel attacked and extend those muscles to lock themselves in. There is some research indicating that roving coral groupers (Plectropomus pessuliferus) sometimes cooperate with giant morays in hunting.
Deepwater benthopelagic fish are robust, muscular swimmers that actively cruise the bottom searching for prey. They often live around features, such as seamounts, which have strong currents. Commercial examples are the orange roughy and Patagonian toothfish. Because these fish were once abundant, and because their robust bodies are good to eat, these fish have been commercially harvested.
Conservation status
Major demersal fishery species in the North Sea such as cod, plaice, monkfish and sole, are listed by the ICES as "outside safe biological limits."
The by-catch problem
A major problem in conservation of demersal fish populations is that of by-catch, whereby fish are caught by accident when targeting other species. The European Commission has written that “A key issue is that many of the most important demersal stocks (i.e. those that live on or near the bottom of the sea) are caught in mixed fisheries. In practice, this means that each time a vessel retrieves its fishing gear, its catch will consist of a mix of different species.” This has led to a situation whereby, even when the International Council for the Exploration of the Sea recommends a Total Allowable Catch of zero for a given demersal species in order to allow replenishment of population, the European Council nonetheless sets the Total Allowable Catch far above zero so long as the catch is by-catch, in order not to prevent trawlers fishing for other species. This means that those threatened species do not get the chance to replenish even when not directly targeted by trawlers.
The Common Sole, Solea solea, is sufficiently broadly distributed that it is not considered a threatened species; however, overfishing in Europe has produced severely diminished populations, with declining catches in many regions. For example, the western English Channel and Irish Sea sole fisheries face potential collapse according to data in the UK Biodiversity Action Plan.
Sole, along with the other major bottom-feeding fish in the North Sea such as cod, monkfish, and plaice, is listed by the ICES as "outside safe biological limits." Moreover, they are growing less quickly now and are rarely older than six years, although they can reach forty. World stocks of large predatory fish and large ground fish such as sole and flounder were estimated in 2003 to be only about 10% of pre-industrial levels. According to the World Wildlife Fund in 2006, "of the nine sole stocks, seven are overfished with the status of the remaining two unknown." Data is insufficient to assess the remaining stocks; however, landings for all stocks are at or near historical lows."
World stocks of large predatory fish and large ground fish such as sole and flounder were estimated in 2003 to be only about 10% of pre-industrial levels, largely due to overfishing. Most overfishing is due to the extensive activities of the fishing industry. Current research indicate that the flounder population could be as low as 15 million due to heavy overfishing and industrial pollution along the Gulf of Mexico surrounding the coast of Texas.
Seafood Watch have placed on their list of seafood that sustainability-minded consumers should avoid the following demersal fish: sturgeon (imported wild), Chilean seabass, cod (Atlantic, imported Pacific), flounder (Atlantic), halibut (Atlantic), sole (Atlantic), grouper, monkfish, orange roughy, demersal shark, red snapper and tilapia (Asia farmed).
| Biology and health sciences | Fishes by habitat | Animals |
2859751 | https://en.wikipedia.org/wiki/Pouch%20%28marsupial%29 | Pouch (marsupial) | The pouch is a distinguishing feature of female marsupials and monotremes, and rarely in males as well, such as in the yapok and the extinct thylacine. The name marsupial is derived from the Latin marsupium, meaning "pouch". This is due to the occurrence of epipubic bones, a pair of bones projecting forward from the pelvis. Marsupials give birth to a live but relatively undeveloped foetus called a joey. When the joey is born it crawls from inside the mother to the pouch. The pouch is a fold of skin with a single opening that covers the teats. Inside the pouch, the blind offspring attaches itself to one of the mother's teats and remains attached for as long as it takes to grow and develop to a juvenile stage.
Variations
Pouches are different amongst different marsupials, two kinds distinguishable (on the front or belly): opening towards the head and extending the cavity under the skin towards the tail (forward, or up) or opening towards the tail and extending towards the front legs (to the rear, backward or down). For example for quolls and Tasmanian devils, the pouch opens to the rear and the joey only has to travel a short distance to get to the opening (resting place) of the pouch. While in the pouch they are permanently attached to the teat and once the young have developed they leave the pouch. The kangaroo's pouch opens horizontally on the front of the body, and the joey must climb a relatively long way to reach it. Kangaroos and wallabies allow their young to live in the pouch well after they are physically capable of leaving, often keeping two joeys in the pouch, one tiny and one fully developed. In kangaroos, wallabies and opossums, the pouch opens forward or up.
Female koalas have been described as having a ‘backward-opening’ pouch like wombats, as opposed to an upward-opening pouch like kangaroos, but that is not true. When a female koala gives birth to young her pouch opening faces neither up nor down, although it is located towards the bottom of the pouch rather than at the top. It faces straight outwards rather than ‘backwards’. It sometimes appears to be ‘backward-facing’ because when the joey is older and leans out of the pouch, this pulls the pouch downwards or ‘backwards’. The pouch has a strong sphincter muscle at the opening to prevent the joey from falling out.
In wombats and marsupial moles, the pouch opens backward or down. Backwards facing pouches would not work well in kangaroos or opossums as their young would readily fall out. Similarly, forward-facing pouches would not work well for wombats and marsupial moles as they both dig extensively underground. Their pouches would fill up with dirt and suffocate the developing young. Kangaroo mothers will lick their pouches clean before the joey crawls inside. Kangaroo pouches are sticky to support their young joey. Koalas are unable to clean out their pouches since they face backwards, so just prior to giving birth to the young koala joey, a self-cleaning system is activated, secreting droplets of an anti-microbial liquid that cleans it out. In a relatively short time, the cleansing droplets clean out all of the crusty material left inside, leaving an almost sterile environment ready to receive the tiny joey.
The water opossum and the now extinct Tasmanian tiger are the only two marsupials where the male also has a pouch (in order to protect their genitalia while swimming).
Some marsupials (e.g. phascogales) lack the true, permanent pouches seen in other species. Instead, they form temporary skin folds (sometimes called "pseudo-pouches") in the mammary region when reproducing. This type of pouch also occurs in echidnas which are monotremes.
Pouches have their own microbiota and it changes depending on reproductive stage: anoestrus, pre-oestrus, oestrus/birth, post-birth. Changes can be made to the appearance of the pouch by inducing oestrus.
| Biology and health sciences | Reproductive system | Biology |
2863567 | https://en.wikipedia.org/wiki/The%20Iron%20Bridge | The Iron Bridge | The Iron Bridge is a cast iron arch bridge that crosses the River Severn in Shropshire, England. Opened in 1781, it was the first major bridge in the world to be made of cast iron. Its success inspired the widespread use of cast iron as a structural material, and today the bridge is celebrated as a symbol of the Industrial Revolution.
The geography of the deep Ironbridge Gorge, formed by glacial action during the last ice age, meant that there are industrially useful deposits of coal, iron ore, limestone and fire clay present near the surface where they are readily mined, but also that it was difficult to build a bridge across the river at this location. To cope with the instability of the banks and the need to maintain a navigable channel in the river, a single span iron bridge was proposed by Thomas Farnolls Pritchard. After initial uncertainty about the use of iron, construction took place over 2 years, with Abraham Darby III responsible for the ironwork. The bridge crosses the Ironbridge Gorge with a main span of , allowing sufficient clearance for boats to pass underneath.
In 1934 it was designated a scheduled monument and closed to vehicular traffic. Tolls for pedestrians were collected until 1950, when the bridge was transferred into public ownership. After being in a poor state of repair for much of its life, extensive restoration works in the latter half of the 20th century have protected the bridge. The bridge, the adjacent settlement of Ironbridge and the Ironbridge Gorge form the UNESCO Ironbridge Gorge World Heritage Site.
History
Background
The Ironbridge Gorge was formed at the end of the last ice age by the overflowing of Lake Lapworth, which resulted in the exposure of useful deposits of resources such as coal, iron ore, fire clay and limestone near the surface where they were readily mined. With the river providing a means of transport, the local area was an important centre of the emerging Industrial Revolution.
Abraham Darby I first smelted local iron ore with coke made from Coalbrookdale coal in 1709, and in the coming decades Shropshire became a centre for industry due to the low price of fuel from local mines. The River Severn was used as a key trading route, but it was also a barrier to travel around the deep Ironbridge Gorge, especially between the then important industrial parishes of Broseley and Madeley, the nearest bridge being at Buildwas away. The Iron Bridge was therefore proposed to link the industrial town of Broseley with the smaller mining town of Madeley and the industrial centre of Coalbrookdale. The use of the river by boat traffic and the steep sides of the gorge meant that any bridge should ideally be of a single span, and sufficiently high to allow tall ships to pass underneath. The steepness and instability of the banks was problematic for building a bridge, and there was no point where roads on opposite sides of the river converged.
The Iron Bridge was the first of its kind to be constructed, although not the first to be considered nor the first iron bridge of any kind. An iron bridge was partly constructed at Lyons in 1755, but was abandoned for reasons of cost, and a span wrought iron footbridge over an ornamental waterway was erected in Kirklees, Yorkshire, in 1769.
Proposal
In 1773, architect Thomas Farnolls Pritchard wrote to his 'iron mad' friend and local ironmaster, John Wilkinson of Broseley, to suggest building a bridge out of cast iron. Although he specialised in the design of chimneypieces and other items of interior decoration, and in funerary monuments, he had also previously designed both wooden and stone bridges.
During the winter of 1773–74, local newspapers advertised a proposal to petition Parliament for leave to construct an iron bridge with a single span. In 1775, a subscription raised funds of between £3000 and £4000 (equivalent to £ to £ in 2016), and Abraham Darby III, the grandson of Abraham Darby I and an ironmaster working at Coalbrookdale, was appointed treasurer to the project.
In March 1776, the Act to build a bridge received Royal assent. It had been drafted by Thomas Addenbrooke, secretary of the trustees, and John Harries, a London barrister, then presented to the House of Commons by Charles Baldwyn, MP for Shropshire. Abraham Darby III was commissioned to cast and build the bridge. In May 1776, the trustees withdrew Darby's commission, and instead advertised for plans for a single arch bridge to be built in "stone, brick or timber". No satisfactory proposal was made, and the trustees agreed to proceed with Pritchard's design, but there was continued uncertainty about the use of iron, and conditions were set on the cost and duration of the construction. In July 1777 the span of the bridge was decreased to , and then increased again to , possibly in order to accommodate a towpath.
Construction
The site, adjacent to where a ferry had run between Madeley and Benthall, was chosen for its high approaches on each side and the relative solidity of the ground. The Act of Parliament described how the bridge was to be built from a point in Benthall parish near the house of Samuel Barnett to a point on the opposite shore near the house of Thomas Crumpton. Pritchard died on 21 December 1777 in his towerhouse at Eyton on Severn, only a month after work had begun, having been ill for over a year.
The bridge is built from five sectional cast-iron ribs that give a span of . The construction of the bridge used of iron, and there are almost 1,700 individual components, the heaviest weighing . Components were cast individually to fit with each other, rather than being of standard sizes, with discrepancies of up to several centimetres between 'identical' components in different locations.
The masonry and abutments were constructed between 1777 and 1778, and the ribs were lifted into place in the summer of 1779. The bridge first spanned the river on 2 July 1779, and it was opened to traffic on 1 January 1781.
In 1997, a watercolour by Elias Martin was discovered in a Stockholm museum, which showed the bridge under construction in 1779. The painting shows a moveable wooden scaffold consisting of derrick poles standing in the river bed being used as a crane to position the half-ribs of the bridge, which had been taken to the site by boat from Darby's foundry downstream. Using the approach depicted in the painting, a half-size replica of the main section of the bridge was built in 2001 as part of the research for the BBC's Timewatch programme, which was shown the following year.
Design
The bridge is to a carpenters' design typically used for wood structures, built from five sectional cast-iron ribs that give a span of . The interlocking rings between the arches make this a truss arch. Exactly of iron was used in the construction of the bridge, and there are almost 1,700 individual components, the heaviest weighing . Components were cast individually to fit with each other, rather than being of standard sizes, with discrepancies of up to several centimetres between 'identical' components in different locations.
Decorative rings and ogees between the structural ribs of the bridge suggest that the final design was Pritchard's, as the same elements appear in a gazebo he rebuilt. A foreman at the foundry, Thomas Gregory, drew the detailed designs for the members, resulting in the use of carpentry jointing details such as mortise and tenon joints and dovetails.
The two outer ribs are engraved with the words: "This bridge was cast at Coalbrook-Dale and erected in the year MDCCLXXIX" (Roman numerals for 1779).
Two supplemental arches, of similar cast iron construction, carry a towpath on the southern bank and also act as flood arches. A stone arch with a brick vault carries a small path on the northern (town side) bank.
Material
The Iron Bridge is made of cast iron. This is immensely strong in compression, but performs less well than steel and wrought iron when subjected to tension or bending moments, because of its brittleness and lower tensile strength. Analysis of an arch and a strut from the Iron Bridge revealed the following elemental compositions:
The presence of 0.1% sulphur in cast iron is at the upper limit of what is acceptable, but the presence of sufficient manganese leads to the formation of harmless manganese sulphide.
Puddled wrought iron was a much better structural material than cast iron but not widely available until after 1800, eventually becoming the preferred material for bridges, rails, ships and buildings until new steel making processes such as the Bessemer process were developed in the late 19th century.
Cost
Darby had agreed to construct the bridge with a budget of £3,250 (equivalent to £426,353 in 2023) and this was raised by subscribers to the project, mostly from Broseley. While the actual cost of the bridge is unknown, contemporary records suggest it was as high as £6,000 (£787,113 in 2023), and Darby, who was already indebted from other ventures, agreed to cover the excess. However, by the mid-1790s the bridge was highly profitable, and tolls were giving the shareholders an annual dividend of 8 per cent.
Later history
The opening of the bridge resulted in changes in the pattern of settlement in the gorge, and roads around the bridge were improved in the years after its construction. The town of Ironbridge, taking its name from the bridge, developed at the northern end. The trustees, as well as local hotel keepers and coach operators, promoted interest in the bridge among members of high society.
Repairs
In July 1783, a wall was built in order to prevent the north bank from slipping into the river. Cracks were found in the stone land arch on the south side in December 1784, and the neighbouring abutment showed signs of movement. The Gorge is very prone to landslides, and over 20 are recorded in the British Geological Survey's National Landslide Database in the area. It was suspected that the sides of the gorge were moving towards the river, forcing the feet of the arch towards each other, and consequently repairs were carried out in 1784, 1791 and 1792.
It was the only bridge on the River Severn to survive the flood of February 1795 undamaged, due to its strength and small profile against the floodwaters. The medieval bridge at Buildwas was replaced with a cast iron bridge by Thomas Telford, which, by virtue of superior design, required half the quantity of iron despite a longer span of . The Buildwas Bridge survived until 1906. Contrary to a common misconception, Telford was not responsible for the construction of the bridge itself.
In 1800 the trustees commissioned repairs which lasted for several years, which involved the replacement of the stone land arches with wooden ones to relieve pressure on the main span. A proposal to build a rigid support between the abutments to keep them apart was found to be impossible with the available technology, but was achieved during the later restoration of the bridge in the 1970s. In 1812, its construction was described as "very bad" by Charles Hutton, and he predicted that it would not last for long, "though not from any deficiency in the iron-work", but due to cracks that had appeared in the stonework. The timber arches were replaced with cast iron ones in December 1820, and further repairs were necessary throughout the remainder of the 19th century.
Around 1870, "Sir B. Baker stated that it had required patching for ninety years, because the arch and the high side arches would not work together. Expansion and contraction broke the high arch and the connexions between the arches. When it broke they fished it. Then the bolts sheared or the ironwork broke in a new place. He advised that there was nothing unsafe; it was perfectly strong and the stress in vital parts moderate. All that needed to be done was to fish the fractured ribs of the high arches, put oval holes in the fishes, and not screw up the bolts too tight."
On 24 August 1902, a length of parapet collapsed into the river, and a section of deck plate weighing around fell from the bridge in July 1903. The opening of a toll-free concrete bridge in 1909 caused concern among the trustees, but it continued to be used by vehicles and pedestrians.
Closure to vehicles
A 1923 report by engineering consultants Mott, Hay and Anderson suggested that other than the paintwork, the main span of the bridge was in good condition. It was suggested that the metal deck of the bridge was dangerously heavy, and that after removing the dead weight the bridge should be reopened to vehicles no heavier than 2 tons and restricted to the centre of the roadway. A weight limit of 4 tons was imposed, but the housing boom of the 1930s meant that drivers distributing tiles produced at Jackfield were insistent that they should be allowed to use the bridge, so the trustees took the decision to close it to vehicular traffic with effect from 18 June 1934. Tolls for pedestrians were collected until 1950, when ownership of the bridge was transferred to Shropshire County Council. The tolls collected only marginally covered the cost of collection, leaving no budget for conservation, and the bridge had not been cleaned or painted for many years. Due to its poor condition, between the 1940s and 1970s a number of suggestions were made to scrap and replace the bridge, or move it to a different location. In 1956 the County Council made a proposal to demolish the bridge and replace it with a new one, but this plan did not come to fruition. The Ironbridge Gorge Museum Trust was set up in 1967 with the aim of protecting industrial heritage in the Ironbridge Gorge, and was able to secure funding from the council to carry out repairs.
Restoration
With funding from Shropshire County Council, the Historic Buildings Council for England, and the nascent Ironbridge Gorge Museum Trust, a programme of repairs took place on the foundations of the bridge at a cost of £147,000 between 1972 and 1975. The consulting engineers Sandford, Fawcett, Wilton and Bell decided to place a ferro-concrete inverted arch under the river to counter inward movement of the bridge abutments. The arch was built by the Tarmac Construction Company, starting in the spring of 1973, but unusually high summer floods washed over the cofferdam, frustrating hopes that the work could be done in a single summer. Filling material was removed from the south abutment to reduce its weight, and the arch through it was reinforced with concrete. The road surface was replaced with a lighter tarmac, the stone of the abutments was renewed and the toll-house was restored as an information centre. In 1980, the structure was painted for the first time in the 20th century, and the work was complete for the bicentenary of the opening, which was celebrated on 1 January 1981.
Between 1999 and 2000, the bridge was scaffolded to allow examination by English Heritage. The bridge was also repainted and minor repairs were carried out. In January 2017 English Heritage announced a £1.2 million restoration project on the Iron Bridge, starting in September 2017, the "biggest ever conservation project" undertaken by English Heritage. The cost was quoted in 2018 at £3.6 million, with English Heritage describing it as "an ambitious conservation of its ribs and arches, its stonework and decking." The project was created after extensive surveys of the area revealed that the historic structure was under threat due to stresses in the ironwork dating from the original construction, ground movement over the centuries, and an earthquake in the 19th century. Apart from the structural restoration, the bridge was also reverted from blue-grey to its original red-brown colour after forensic analysis revealed this was how the bridge looked when it was first erected.
The project was partly funded through the use of crowdfunding in which £47,545 was raised. It also received a €1,000,000 donation from the Hermann Reemtsma Foundation, a German foundation which mainly promotes cultural and social projects in northern Germany. This would also be the foundation's first funding in the United Kingdom. Following the works the bridge was reopened on 6 December 2018.
Recognition
The bridge, the adjacent settlement of Ironbridge and the Ironbridge Gorge form the UNESCO Ironbridge Gorge World Heritage Site, which was created in 1986. The bridge is a Grade I listed building, and is owned by Telford and Wrekin Council.
In 1934 it was designated a Scheduled Ancient Monument, and in 1979, the bridge was recognised by the American Society of Civil Engineers as an International Civil Engineering Landmark.
In 2020 the conservation and restoration work on the bridge was honoured in the European Heritage Awards/Europa Nostra Awards, which includes up to 30 of the most outstanding heritage projects from all parts of Europe within the EU and EEA.
Influence on bridge design
The bridge, as the first bridge of significant size built of metal, had "considerable influence on developments in the fields of technology and architecture".
The successful use of cast iron in 1781 pioneered the choice of that material for many subsequent bridges, and cast iron arches of considerable span were constructed late in the 18th and early in the 19th century.
In 1786, the revolutionary author and polymath Thomas Paine had models built to demonstrate the use of cast iron for bridges, and promoted these to the Academy of Sciences at Paris, and to the Royal Society in England. He went on to have manufactured and erected "a complete rib of 90 feet span, and 5 feet of height from the chord line to the center of the arch", weighing three tons. This was followed by a complete bridge of five ribs and span which he had erected in a field in Paddington, but without buttresses it was merely for display and was dismantled after a year.
In 1793–96, the Wearmouth Bridge was built with a span of , constructed from cast iron in the form of cast voussoirs, somewhat like the voussoirs of a masonry bridge. The bridge used some of the iron from Paine's bridge, which had been returned to the foundry in Rotherham.
In 1795, a large flood in the Severn swept away all the bridges in the vicinity, except the Iron Bridge, where the open structure allowed the floodwaters to pass through. Thomas Telford was Surveyor of Public Works in Shropshire at the time. His design for the replacement bridge at Buildwas incorporated a high arch like the Iron Bridge, but it had a span wider and used less than half the amount of iron. Telford went on to design a series of cast iron bridges, the oldest of which to survive is the Craigellachie Bridge.
In 1799, the Coalport Bridge was rebuilt after the flood as a single span with three cast iron ribs. This 1799 version was re-modelled in 1818 with two additional ribs, and survives to the present day.
Artistic depictions
Over fifty painters and engravers came to the area around Coalbrookdale between 1750 and 1830 to witness and record the rise of industry and changing landscape. One of the first artists to depict the bridge was William Williams, who was paid 10 guineas () in October 1780 by Darby for a drawing of the bridge. An engraving by Michael Angelo Rooker proved popular, and a copy was purchased by Thomas Jefferson where it was displayed in the dining room of Monticello.
In 1979, the Royal Academy of Arts held an exhibition entitled "A View from the Iron Bridge" to commemorate the bicentenary of the bridge.
| Technology | Bridges | null |
3833397 | https://en.wikipedia.org/wiki/Influenza%20B%20virus | Influenza B virus | Influenza B virus is the only species in the genus Betainfluenzavirus in the virus family Orthomyxoviridae.
Influenza B virus is only known to infect certain mammal species, including humans, ferrets, pigs, and seals. This limited host range is apparently responsible for the lack of influenza pandemics associated with influenza B virus, in contrast with those caused by the morphologically similar influenza A virus, as both mutate by both antigenic drift and reassortment. Nevertheless, it is accepted that influenza B virus could cause significant morbidity and mortality worldwide, and significantly impacts adolescents and schoolchildren.
There are two known circulating lineages of influenza B virus based on the antigenic properties of the surface glycoprotein hemagglutinin. The lineages are termed B/Yamagata/16/88-like and B/Victoria/2/87-like viruses. The quadrivalent influenza vaccine licensed by the CDC has been designed to protect against both co-circulating lineages and as of 2016 has been shown to have greater effectiveness in prevention of influenza caused by influenza B virus than the previous trivalent vaccine.
However, the B/Yamagata lineage might have become extinct in 2020/2021 due to COVID-19 pandemic measures. In October 2023, the World Health Organization concluded that protection against the Yamagata lineage was no longer necessary in the seasonal flu vaccine, reducing the number of lineages targeted by the vaccine from four to three. For the 2024–2025 Northern Hemisphere influenza season, the US Food and Drug Administration (FDA) recommends removing B/Yamagata from all influenza vaccines. The European Medicines Agency (EMA) recommends removing B/Yamagata from influenza vaccines for the 2024–2025 seasonal flu vaccine composition.
Morphology
The influenza B virus capsid is enveloped while its virion consists of an envelope, a matrix protein, a nucleoprotein complex, a nucleocapsid, and a polymerase complex. It is sometimes spherical and sometimes filamentous. Its 500 or so surface projections are made of hemagglutinin and neuraminidase.
Genome structure and genetics
The influenza B virus genome is 14,548 nucleotides long and consists of eight segments of linear negative-sense, single-stranded RNA. The multipartite genome is encapsidated, each segment in a separate nucleocapsid, and the nucleocapsids are surrounded by one envelope.
The ancestor of influenza viruses A and B and the ancestor of influenza virus C are estimated to have diverged from a common ancestor around 8,000 years ago. Influenza viruses A and B are estimated to have diverged from a single ancestor around 4,000 years ago, while the subtypes of influenza A virus are estimated to have diverged 2,000 years ago. Metatranscriptomics studies have also identified closely related "influenza B-like" viruses such as the Wuhan spiny eel influenza virus and also "influenza B-like" viruses in a number of vertebrate species such as salamanders and fish.
Diminishing the impact of this virus is the fact that, "in humans, influenza B viruses evolve slower than A viruses and faster than C viruses". Influenza B virus mutates at a rate 2 to 3 times slower than type A.
Vaccine
In 1936, Thomas Francis Jr. discovered the ferret influenza B virus. Also in 1936, Macfarlane Burnet made the discovery that influenza virus may be cultured in hen embryonated eggs. This prompted research into the properties of the virus and the creation and application of inactivated vaccines in the late 1930s and early 1940s. Inactivated vaccines' usefulness as a preventative measure was proven in the 1950s. Later, 2003 saw the approval of the first live, attenuated influenza vaccine. Looking into influenza B specifically, Thomas Francis Jr. isolated influenza B virus in 1936. However, it was not until 1940 that influenza B viruses were discovered.
In 1942, a new bivalent vaccine was developed that protected against both the H1N1 strain of influenza A and the newly discovered influenza B virus. In today's current world, even while some technology has advanced and flu vaccines now cover both strains of influenza A and B, the science is still based on findings from almost a century ago. The viruses included in flu vaccines are changed each year to match the strains of flu that are most likely to make people sick that year since flu viruses can develop swiftly and new mutations have appeared each year, like H1N1.
Even though there have been two different lineages of influenza B viruses that were circulating during most seasons, flu vaccinations were long meant to protect against three different flu viruses: the influenza A(H1N1), influenza A(H3N2), and one type of influenza B virus. The second lineage of the B virus was since added to provide greater defense against circulating flu viruses. Two influenza A viruses and two influenza B viruses have up until 2023 been among the four flu viruses that a quadrivalent vaccine was intended to protect against. As of 2022 all flu vaccines in the United States were quadrivalent. The four main types of type A and B influenza viruses that are most likely to spread and make people sick during the upcoming flu season have been the targets of seasonal influenza (flu) vaccines. All of the available flu vaccinations in the United States have offered protection against the influenza A(H1), A(H3), B/Yamagata, and B/Victoria lineage viruses. Each of these four vaccine virus components has been chosen based on which flu viruses are infecting people ahead of the upcoming flu season, how widely they are spreading, how well the vaccines from the previous flu season may protect against those flu viruses, and the vaccine viruses' capacity to offer cross-protection.
For the 2022–2023 flu season, there were three flu vaccines that were preferentially recommended for people 65 years and older; various influenza (flu) vaccinations are authorized for use in people of various age groups. In March 2022, the Vaccines and Related Biological Products Advisory Committee of the US Food and Drug Administration (FDA) proposed using A(H1N1)pdm09, A(H3N2), and B/Austria/1359417/2021-like viruses for trivalent influenza vaccines to be utilized in the US.
However, the B/Yamagata lineage might have become extinct in 2020/2021 due to COVID-19 pandemic measures, and there have been no naturally occurring cases confirmed since March 2020. In October 2023, the World Health Organization concluded that protection against the Yamagata lineage was no longer necessary in the seasonal flu vaccine, reducing the number of lineages targeted by the vaccine from four to three. If the virus has indeed gone extinct, it would be the first documented case of a virus going extinct due to changes in human behavior. However, the overall burden of the influenza virus would be similar to past seasons because there are still three remaining strains widely circulating.
Discovery and development
In 1940, an acute respiratory illness outbreak in Northern America led to the discovery of influenza B virus (IBV), which was later discovered to not have any antigenic cross-reactivity with influenza A virus (IAV). Based on calculations of the rate of amino acid substitutions in HA proteins, it was estimated that IBV and IAV diverged from one another around 4000 years ago. However, the mechanisms of replication and transcription, as well as the functionality of the majority of viral proteins, appear to be largely conserved, with some unusual differences. Although IBV has occasionally been found in seals and pigs, its primary host species is the human. IBVs can also spread epidemics throughout the world, but they receive less attention than IAVs do due to their less prevalent nature, both in infecting hosts and in the symptoms that result from infection. IBVs used to be unclassified, but since the 1980s, they have been divided into the B/Yamagata and B/Victoria lineages. IBVs have further divisions known as clades and sub-clades, just like IAVs do.
Hemagglutinin (HA) and neuraminidase (NA) are two virus surface antigens that are constantly changing. Antigenic drift or antigenic shift are two possible influenza viral changes. Small changes in the HA and NA of influenza viruses caused by antigenic drift result in the creation of novel strains that the immune system of humans might not be able to identify. These emerging strains are the influenza virus's evolutionary responses to a potent immunological response across the population. The main cause of influenza recurrence is antigenic drift, which makes it essential to reevaluate and update the influenza vaccine's ingredient list every year. Annual influenza outbreaks are caused by antigenic drift and declining immunity, when the residual defenses from prior exposures to related viruses are incomplete. Antigenic drift occurs in influenza A, B, and C.
Hemagglutination inhibition experiments using ferret serum after infection allowed the identification of two very different antigenic influenza type B variants in the years 1988–1989. These viruses shared antigens with either B/Yamagata/16/88, a variation that was discovered in Japan in May 1988, or B/Victoria/2/87, the most recent reference strain. The B/Victoria/2/87 virus shared antigens with all influenza B viruses discovered in the United States during an outbreak in the winter of 1988–1989.
In Japan, influenza B virus reinfection was investigated virologically in 1985–1991 and epidemiologically in 1979–1991 in children. Four influenza B virus outbreaks that each included antigenic drift occurred during the course of this study. Between the epidemics in 1987–1988 and 1989–1990, there was a significant genetic and antigenic change in the viruses. Depending on the influenza seasons, the minimum rate of reinfection with influenza B virus for the entire period was between 2 and 25%. Hemagglutination inhibition assays were used to examine the antigens of the influenza B virus primary and reinfection strains that were isolated from 18 children between the years of 1985 and 1990, which encompassed three epidemic periods. The findings revealed that reinfection occurred with the viruses recovered during the 1984–1985 and 1987–1988 influenza seasons, which belonged to the same lineage and were antigenically close.
Today, the B/Yamagata lineage might be extinct as a result of COVID-19 pandemic measures, and there have been no naturally occurring cases confirmed since March 2020. Although this development has resulted in updated recommendations regarding vaccine composition, continued surveillance is required to assess this conclusion fully, as pauses in IBV circulation have been observed before.
| Biology and health sciences | Specific viruses | Health |
3837188 | https://en.wikipedia.org/wiki/Grapsus%20grapsus | Grapsus grapsus | Grapsus grapsus is one of the most common crabs along the western coast of the Americas. It is known as the red rock crab, or, along with other crabs such as Percnon gibbesi, as the Sally Lightfoot crab.
Distribution
Grapsus grapsus is found along the Pacific coast of Mexico, Central America, and South America (as far south as northern Peru), and on nearby islands, including the Galápagos Islands. It is also found along the Atlantic coast of South America, but is replaced in the eastern Atlantic Ocean (Ascension Island and West Africa) by its congener Grapsus adscensionis.
Description
Grapsus grapsus is a typically shaped crab, with five pairs of legs, the front two bearing small, blocky, symmetrical chelae (claws). The other legs are broad and flat, with only the tips touching the substrate. The crab's round, flat carapace is slightly longer than . Young G. grapsus are black or dark brown in colour and are camouflaged well on the black lava coasts of volcanic islands. Adults are quite variable in colour; some are muted brownish-red, some mottled or spotted brown, pink, or yellow.
Taxonomy
Grapsus grapsus was first described by Carl Linnaeus in the 1758 10th edition of Systema Naturae as "Cancer grapsus".
The species Grapsus grapsus and G. adscensionis were not separated until 1990. The latter is found in the eastern Atlantic, while the former is not. While the validity of the separation into two species has been questioned, there are constant morphological differences in the colouration of the pereiopods and the form of the first zoea larva, and no evidence for any genetic connection between the two populations, and they are generally treated as separate species.
Ecology and behavior
This crab lives among the rocks at the often turbulent, windy shore, just above the limit of the sea spray. It feeds on algae primarily, sometimes sampling other plant matter and sponges (such as clams), crustaceans (including other crabs), fishes, young sea turtles, bird eggs and droppings, bat guano and dead animals (mainly seals and birds). As larvae, they feed on phytoplankton. They have been known to resort to cannibalism when populations densities are high or food is scarce. It is an agile crab, capable of leaping, and consequently hard to catch. Not considered very edible by humans, it is used as bait by fishermen. It is preyed upon by the chain moray eel, Echidna catenata, as well as by octopuses.
G. grapsus has been observed in an apparent cleaning symbiosis taking ticks from marine iguanas on the Galápagos Islands.
Grapsus grapsus was collected by Charles Darwin during his voyages on HMS Beagle, and also by the first comprehensive study of the fauna of the Gulf of California, carried out by Ed Ricketts, together with John Steinbeck and others. Steinbeck records:
These little crabs, with brilliant cloisonné carapaces, walk on their tiptoes, They have remarkable eyes and an extremely fast reaction time. In spite of the fact that they swarm on the rocks at the Cape [San Lucas], and to a less degree inside the Gulf [of California], they are exceedingly hard to catch. They seem to be able to run in any of four directions; but more than this, perhaps because of their rapid reaction time, they appear to read the mind of their hunter.
Gallery
| Biology and health sciences | Crabs and hermit crabs | Animals |
3838790 | https://en.wikipedia.org/wiki/Arthropleura | Arthropleura | is an extinct genus of massive myriapod that lived in what is now Europe and North America around 345 to 290 million years ago, from the Viséan stage of the lower Carboniferous Period to the Sakmarian stage of the lower Permian Period. It is related to millipedes, and was capable of reaching at least in length, possibly up to over , making it the largest known land arthropod of all time. Arthropleura is known from body fossils as well as trace fossils, particularly giant trackways up to wide, and potentially also large burrows. It lived in open, sparsely wooded environments near water, and was possibly amphibious.
History and classification
First discovered in 1854, Arthropleura has consistently attracted much artistic and scientific attention, yet has historically been known from mostly fragmentary remains. Prior to its description, the remains were attributed by Jordan and Meyer to a decapod crustacean. In their subsequent description of the genus, they compared it with trilobites and eurypterids. Another author, Moritz Kliver, described a small specimen of Arthropleura armata showing the underside, with a series of 7 pairs of limbs articulated with various other plates, including the sternites. Kliver that Arthropleura could not be an insect, arachnid, or myriapod, but instead that it was a non-decapod crustacean, comparing the appendages to those of branchiopods. Following authors tended to view Arthropleura in connection to isopods, some believing it represented a very primitive Eumalacostracan crusteacean. Palaeontologists Lewis Moysey and Henry Woodward were the first to associate Arthropleura with Myriapoda in 1911.
Gérard Waterlot published a landmark study in 1934, establishing the order Arthropleurida, placing the group as sister to order Trilobita, within the defunct crustacean subclass Archaeocrustacea. He believed the limbs to be composed of two branches, a lower walking leg and an upper gill branch, like the biramous limbs of trilobites. Palaeontologist Leif Størmer remarked on this study in 1944, writing that after careful study of the photographs, the supposed two branches of the limb are always found close together, close enough that even their joints consistently line up. Because of this, he interpreted the limbs as uniramous, having only a single branch. Waterlot referred a nearly complete juvenile specimen described by Dr. W. T. Calman as Arthropleura moyseyi to Arthropleura armata, and agreed that this specimen showed the mostly unknown head segment of the animal. A curved structure found along the edge of the head resembled the mandibles of some myraipods, and so it was seen by Waterlot as a cephalic feeding appendage, and that this meant Arthropleura was carnivorous, feeding on small soft-bodied prey. Waterlot also believed that Arthropleura was amphibious, living at the bottom of lakebeds, with occasional excursions onto the land where the humidity of coal swamps allowed it to continue breathing with its supposedly trilobite-like gills. Other authors preferred a terrestrial habitat, and questioned the carnivorous diet. The juvenile specimen described by Waterlot appeared to show various plant remains within the gut tract, more recently this has been seen as a taphonomic artifact. The most complete known adult specimen of Arthropleura (the Maybach specimen) was described by Paul Guthörl in 1935, discovered in the roof of a German coal mine and carefully extracted, measuring long.
Arthropleura as a myriapod or relative of myriapods (particularly diplopods) becaming the prevailing view among scientists in the following decades, and the limbs became quite well understood. Further species of Arthropleura have been described over the decades, mostly from central Europe and the UK, and even from the central United States with Arthropleura cristata. North American remains are also known from widespread walking traces, as well as fossil remains from Nova Scotia, Ohio, and Pennsylvania. The head segment, however, remained enigmatic even into the 21st century. What had previously been identified as the head was reinterpreted as the collum, the first segment behind the head, with the true head hidden beneath it. The Devonian millipede Microdecemplex was described in 1999, and was believed to belong to the subclass Arthropleuridea. Thse tiny fossils were much more complete, and so its head anatomy was taken as an insight into the head of Arthropleura.
Finally, in 2024, exceptionally preserved fossils from the Montceau-les-Mines Lagerstätte of France were studied and described with the aid of Micro-CT scanning, revealing the head region in two nearly complete juvenile Arthropleura fossils. They showed that the head identified by past authors was indeed the head, with the collum a thin segment behind it. Waterlot's cephalic limbs are now identified as paired ventral scleites, protecting the head like a car bumper. The cephalic limbs showed two pairs of external maxillae, and an internal set of mandibles. Phylogenetic analysis found Arthropleura just outside of Diplopoda (stem-group Diplopoda), while Microdecemplex was found as a crown-group millipede rather than an arthropleuridean.
Morphology
Multiple partial specimens and trackways suggest that Arthropleura could exceed long. Tracks from Arthropleura up to wide have been found at Joggins, Nova Scotia. In 2021 a fossil, probably a shed exoskeleton (exuviae) of an Arthropleura (indeterminate species, as due to its exposure of only the underside of the carapace, ornamentation cannot be observed), was reported with an estimated width of , length of to and body mass of . It is one of the largest arthropods ever known, as large as the eurypterid Jaekelopterus rhenaniae, whose length is estimated at . The 2024 study reported the complete head and trunk of a juvenile specimen of Arthropleura sp. (MNHN.F.SOT002123) from Kasimovian (~305 Ma) Montceau-les-Mines lagerstätte, which revealed multiple previously unknown features. Arthropleura had large, flattened ventral sclerites and a pair of antennae with at least seven antennal articles at the front of its head. The trunk anatomy of Arthropleura is characterized by a series of well-developed 28-32 tergites (dorsal exoskeleton) having three lobes like a trilobite, the dorsal surface of which is typically covered by many tubercles or spines. Juvenile specimens have fewer numbers of tergites at 20–24, suggesting hemianamorphic development, with the number of segments reaching their maximum number (adding during each moult) before the final moult, the animal continuing to grow while retaining the same number of segments past a certain point. Arthropleura would have had pleurites, paired sclerites between the sternite and tergite on either side of the body. They would have been beneath the paratergites and above the legs, held in place by arthrodial membrane. Pleurites have been identified in a number of specimens (such as in the Maybach specimen) as oval-shaped plates with rough granulose ornamentation, and would have provided an attachment point for trunk musculature.
Arthropleura had two pairs of maxillae appendages beneath the head, each composed of 3 segments (podomeres). Their morphology are similar to those of a centipede: the outermost pair (second maxillae) being relatively long and pointed, while the inner pair (first maxillae) were short and plate-like. The paired mandibles were small, composed of three segments, and fully internalized into the head, similar to centipedes rather than millipedes. The segmentation of antennae and mandibles identify it as a millipede, but the presence of a pair of legs underneath its collum, which is absent in present-day species, reveals it as the sister group of the millipede crown group. Unlike any living myriapods, they had stalked compound eyes, which despite also being known from the extinct stem-myriapod group Euthycarcinoidea, appears to be a derived trait evolved independently from the euthycarcinoids. Each body segment bore two pairs of walking legs, which themselves are composed of 9 or 10 segments (podomeres). A crease ran down each side of the leg which probably allowed for stronger muscle attachment as an apodeme, and the ventral surface bore paired endites on each segment (actually macro-setae, set in sockets and presumably mobile). The anterior surface of the leg was smooth, while the posterior surface was tuberculate, with many pores that housed sensory setae.
Around each walking leg pair, there were three pairs of ventral plates located alongside the median sternite, namely K-plates, B-plates, and rosette plates, and either the B-plates or K-plates were thought to be respiratory organs. This has however also been questioned, with Rolfe and Ingham (1967) considering all of these plates to be simply sclerotized ventral integument responsible for reinforcing and buttressing the limb bases to enable locomotion for such a large animal. Wilson (1999) disagreed with this interpretation, stating that the K-plate was actually sac-shaped with a hollow interior. The outer surface of the K-plate was covered in pitted setae, and the interior of the K-plate was reticulate, possibly serving as a respiratory organ. In this case, movement of the limbs would have helped pump hemolymph in and out of the respiratory apparatus. The spiracle opening to the lungs has never been identified, and the respiratory abilities of Arthropleura are still mostly unknown. The rosette plate is divided into small triangular lobes - the creases between them probably functioned as internal apodemes, muscule attachment for the limbs. The B-plate, which may be homologous with the coxa or basal leg segment in other arthropods, is continuous with the rosette plate. The placement of all these plates in front of the limb also suggests the legs would thrust backward and downward during movement, before returning to their standard position, with little capacity for movement forward. The body terminated with a small, tuberculate, trapezoidal telson.
The different species of Arthropleura are primarily differentiated by the dorsal ornamentation of their tergites. There is a keel across the middle of the tergites running laterally in all species, and typically, there is a line of large tubercules running behind it. The only exception is A. cristata from the Mazon Creek of the United States, where tuberculate ornamentation is absent. Because complete specimens are rare, species are sometimes poorly delineated, and could possibly represent different growth stages of the a smaller number of species, leading paleontologists who work with Arthropleura to often only give a genus-level identification to the fossils. The Montceau-les-Mines species (known from only juvenile specimens) has a band of 4 large tubercules on the paratergal lobe, which rise up into tall spines. This pattern is extremely similar to A. moyseyi, which suggests that A. moyseyi is also a juvenile, possibly of A. mammata. Because the ornamentation may change over time, it is not viewed as a reliable way to determine species. At Montceau-les-Mines, large specimens are absent, and the largest trackways were left by animals less than in length, so this species may have been genuinely much smaller than the giant species elsewhere in Europe.
Paleobiology
Most fossils of Arthropleura are believed to represent exuviae (moulted shells) instead of carcasses. Arthropleura has typically been reconstructed living in coal swamps, based on the co-occurrence of its fossils with dense plant remains and coal veins. However, this view is no longer strictly supported. Many fossils of Arthropleura are found either before or after the dominance of such coal forests, and direct geological and paleobotanical evidence suggests its primary habitat was far more open. Arthropleura preferred open environments, particularly sparsely wooded and near to water, such as coastlines and floodplains. Arthropleura is also found in association with tetrapod trackways suggesting they lived in the same environments. In addition to large trackways, enormous burrows have been attributed to Arthropleura and arthropleurids generally, suggested to be made during periods of aestivation similar to hibernation. Arthropleura remains (body fossils and trackways) are also closely associated with the palaeoequator, with all documented fossils occurring with 10° of either side of the palaeoequator during both the Carboniferous and Permian.
Arthropleura has repeatedly been suggested to have been amphibious to some degree. Despite falling out of favour for many years, this idea has been bolstered by new evidence. Trackways, often found in association with coastlines, are found in both submerged and emergent substrates, showing that Arthropleura was probably capable of walking in air and in shallow water, and did so accordingly. Additional evidence for this comes from the juvenile Montceau fossils, which demonstrate that the eyes were large and stalked, unlike any living myriapod. This condition is however known in the extinct stem-myriapod group euthycarcinoidea (with which Arthropleura sometimes co-occurs), which were aquatic to amphibious. The authors suggest that this could point to a semi-aquatic, amphibious lifestyle, capable of entering and exiting shallow bodies of water. It has been suggested that due to their large size, moulting (a stressful period even for smaller arthropods) probably occurred underwater, allowing its weight to be supported while its new exoskeleton hardened. Wilson favoured a more terrestrial habit for Arthropleura, with no obvious mechanism for the animal to breathe underwater. She suggested that during moulting, Arthropleura could have behaved like modern Scolopendra centipedes, moving in a worm-like fashion without use of their softened legs until they were hardened. She suggested the moulting could have been understaken inside constructed burrows, although remained open to Rolfe (1985)'s theory that moulting was done inside hollow lycopod tree trunks, aided by their naturally high humidity which could have prevented the animal from rapidly drying out.
The diet of Arthropleura has also been heavily debated. Waterlot suggested it was a carnivore, but the holotype of "A. moyseyi" was suggested to preserve the original gut contents, a mix of woody plant material, suggesting it was an herbivore. The interpretation of these fossils as gut contents is no longer supported. Currently, Arthropleura is believed to be a detritivore, like most extant millipedes. This means feeding on either dead and decaying plant matter or animal remains when available. The short, closely packed legs, as well as evidence from the morphology of the ichnofossil trackways, both suggest that Arthropleura was a very slow moving animal, and the lack of venomous forcipules or other predatory adaptions to the limbs basically precludes a predatory lifestyle.
The giant size of Arthropleura has been frequently attributed to higher oxygen levels during the Carboniferous. However, this does not align with the fossil record or modern understanding of arthropod size range. Arthropleura reaches giant sizes before the rise in oxygen concentration during the Carboniferous, with the largest known body fossil found in an interval where oxygen levels were only about 23% higher than modern day. More likely, Arthropleura was capable of gigantism due to other generally favourable environmental and ecological conditions and lack of serious competition. Consequently, Arthropleura went extinct during the early Permian, probably due to the desertification of the equatorial regions of the supercontinent and competition with Permian tetrapods.
| Biology and health sciences | Fossil arthropods | Animals |
3839250 | https://en.wikipedia.org/wiki/Common%20bream | Common bream | The common bream (Abramis brama), also known as the freshwater bream, bream, bronze bream, carp bream or sweaty bream, is a European species of freshwater fish in the family Cyprinidae. It is now considered to be the only species in the genus Abramis.
Range and habitat
The common bream's home range is Europe north of the Alps and Pyrenees, as well as the Balkans. They are found as far east as the Caspian Sea, the Black Sea, and the Aral Sea. The common bream lives in ponds, lakes, canals, and slow-flowing rivers.
Description
The bream is usually long, though some specimens of have been recorded; it usually weighs . Its maximum length is , the record weight exceeds .
The common bream has a laterally flattened and high-backed body and a slightly undershot mouth. It has a bright silver colouration, though older fish can be bronze-coloured, especially in clear waters. The fins are greyish to black, but never reddish.
Similar-looking fish
The common bream can easily be confused with the silver or white bream (Blicca bjoerkna), in particular at the younger stages (see picture). The most reliable method of distinguishing these species is by counting the scales in a straight line downwards from the first ray of the dorsal fin to the lateral line. Silver bream have fewer than 10 rows of scales, while common bream have 11 or more. At the adult stage the reddish tint of the fin of the silver bream is diagnostic. Like other Cyprinidae, common bream can easily hybridise with other species, and hybrids with roach (Rutilus rutilus) can be very difficult to distinguish from pure-bred bream.
Immature specimens could also be confused with other European breams, such as the two Ballerus species or Vimba vimba.
Habitat
The common bream generally lives in rivers (especially in the lower reaches) and in nutrient-rich lakes and ponds with muddy bottoms and plenty of algae. It can also be found in brackish sea waters.
Feeding habits
The common bream lives in schools near the bottom. At night, common bream can feed close to the shore, and in clear waters with sandy bottoms, feeding pits can be seen during daytime. The fish's protractile mouth helps it dig for chironomid larvae, Tubifex worms, bivalves, and gastropods. The bream eats water plants and plankton, as well.
In very turbid waters, common bream can occur in large numbers, which may result in a shortage of bottom-living prey such as chironomids. The bream are then forced to live by filter feeding with their gill rakers, Daphnia water fleas being the main prey. As the fish grows, the gill rakers become too far apart to catch small prey and the bream will not then grow bigger than . If a common bream is malnourished, it can develop a so-called "knife back", a sharp edge along its back.
Spawning
The common bream spawns from April to June, when water temperatures are around . At this time, the males form territories within which the females lay 100,000 to 300,000 eggs on water plants. The fry hatch after three to 12 days and attach themselves to water plants with special adhesive glands, until their yolk is used up.
Because of their slender shape, the young fish are often not recognised as bream, but they can be identified by their flat bodies and silvery colour. At this stage, the fish are still pelagic, but after a few months, they acquire their typical body shape and become bottom-dwellers. By three to four years old, the fish are sexually mature.
Fishing
The freshwater bream is not generally caught for consumption. Common bream are popular with sport and match fishermen. However, bream are not as hard fighting as most other fish native to the UK, as due to their flat, disc-shaped profile they are relatively easy to bring to the bank.
Bream will eat most baits, especially:
Sweetcorn – two or three grains hooked or hair-rigged.
Maggots/worms – two or three straight on the hook.
Boilies – the large mouths of Bream will devour most boilies
Bream can be caught in rivers or lakes, with generous use of groundbait to attract the shoals. They are not shy fish. Another technique is float fishing on the bottom. Ledgering (using just a lead weight to hold the bait down) with a cage feeder full of bait often works better on larger rivers and lakes.
the current European record common bream caught with rod and reel is , caught in the United Kingdom.
| Biology and health sciences | Cypriniformes | Animals |
24412554 | https://en.wikipedia.org/wiki/Anoplotheriidae | Anoplotheriidae | Anoplotheriidae is an extinct family of artiodactyl ungulates. They were endemic to Europe during the Eocene and Oligocene epochs about 44—30 million years ago. Its name is derived from the ("unarmed") and θήριον ("beast"), translating as "unarmed beast".
Ecology
Species of Anoplotheriidae varied substantially in size. Diplobune minor is suggested to have weighted about , while Anoplotherium is suggested to have been up to in weight. Anoplotherium is thought to have been a browser that reared up on its hind legs to feed, while Diplobune is suggested to have been an arboreal climbing animal.
Systematics and taxonomy
The family Anoplotheriidae was assigned to Belluae by Bonaparte (who named it Anoplotheriina) in 1850; to Artiodactyla by Cope in 1889, to Ruminantia by Gregory in 1910, and finally to its own superfamily Anoplotherioidea by Romer in 1966. A 2019 study considered them to be closely related to Cainotheriidae, another group of endemic European artiodactyls, with this group in turn being related to ruminants, while a 2020 study found them to be more closely related to the also European endemic Xiphodontidae, again as relatives of ruminants.
Included genera:
Subfamily Anoplotheriinae
Anoplotherium Cuvier, 1804
Diplobune Rutimeyer, 1862
Duerotherium Cuesta & Badiola, 2009
Ephelcomenus Hurzeler, 1938
Robiatherium Sudre, 1988
Subfamily Dacrytheriinae
Dacrytherium Filhol, 1876
Catodontherium Depéret, 1908
| Biology and health sciences | Other artiodactyla | Animals |
1417480 | https://en.wikipedia.org/wiki/Rome%20Metro | Rome Metro | The Rome Metro () is a rapid transit system that operates in Rome, Italy. It started operation in 1955, making it the oldest in the country.
The Metro comprises three lines – A (orange), B (blue) and C (green) – which operate on of route, serving 73 stations. It has a daily ridership of approximately 820,000 passengers, and an annual traffic of approximately 320 million passengers.
In addition to the Metro, the center of Rome and its urban area are served by 8 FL lines (672 km (417.5 mi) with 131 stations) that surround Rome and the Lazio region, 6 tram lines (36 km (22 mi) ) with 192 stations), 3 commuter urban lines (135 km (83.8 mi) with 57 stations), as well as the Leonardo Express which connects Roma Termini, the central station of the city of Rome, to the Leonardo da Vinci Airport of Fiumicino, and the Civitavecchia Express which connects the city to the main port of Rome, the Port of Civitavecchia. Network extensions are currently under construction on Line C (Porta Metronia, Colosseo-Fori Imperiali and Venezia). There are further projects for Line A, Line B, Rome-Giardinetti and for the suburban rail system. The entire transport system in Rome uses the Metrebus integrated tariff system (an acronym composed of the words "Metro", "Train" and "Bus"), which can be used within the limits of the Municipality of Rome and within the limits of the urban tariff.
Line B was the first metro line inaugurated in the system, and the first official metro in Italy, but the names 'A' and 'B' were only added when the second line opened 25 years after the first. Inaugurated in post-war Italy in 1955 during the reconstruction and on the verge of the Italian economic miracle, it was designed and built for the 1942 universal exhibition (Esposizione Universale Roma, which is now the current business center of Rome) desired by the fascist regime, which never took place due to the outbreak of the World War II.
Lines
Line A
Line A runs from the southeastern suburbs of Rome, then along the northeast section of downtown, and then to the northern section of the city, near Vatican City. It connects with Line B, along with many other national and regional rail services, at Termini, and with Line C at San Giovanni. It has 27 stations, with terminals at Battistini and Anagnina. It is identified by the colour orange.
Line A was the second line built in Rome. Approval was given for the construction of the city's second Metro line in 1959.
Work on Line A began in 1964 in the Tuscolana area, but suffered a series of delays caused the originally planned cut and cover method of construction posed serious problems for road traffic in southeast Rome. Work on the Metro was suspended and began again five years later, using bored tunnels, which partially resolved the traffic problems but caused numerous claims for compensation for vibrations caused by the machines. Work was also frequently interrupted by archaeological finds made during the excavations, particularly near Piazza della Repubblica.
Line A entered service in February 1980. In the late 1990s, it was extended from Ottaviano, in the Prati district, to Battistini to the west.
Since June 2022, the station of Valle Aurelia, on Line A, is connected with the reactivated railway station of Vigna Clara. The Vigna Clara-Valle Aurelia section is a relevant step to close the railway ring in North Rome because of the connection with the Line A, Line B (Ostiense) and the FL3 suburban line to Viterbo.
Line B
Line B was the first Metro line in Rome. Line B connects the northeast of the city with the southwest. It has 26 stations with terminals at Rebibbia, Jonio and Laurentina (just east of EUR). It is identified by the colour blue. Transfers are available with Line A and other rail services at Termini station.
Line B was planned during the 1930s by the fascist government to provide a rapid connection between the main train station, Termini, and a new district to the southeast of the city, E42, the planned location of the Universal Exposition (or Expo), which was to be held in Rome in 1942. The exposition never took place due to Italy's entry into the Second World War in 1940. When its construction was interrupted, some of the tunnels on the city-centre side of the Metro (between Termini and Piramide) had already been completed, and they were used as air raid shelters during the war.
Work restarted in 1948, together with the development of the site formerly designated for the Expo into a residential and business district under the name EUR. The Metro was officially opened on 9 February 1955 by the then President of the Republic, Luigi Einaudi. Regular services began the following day.
In 1990, Line B was extended from Termini to Rebibbia in the east of the city and the entire line was modernised. A new long branch (B1) was opened connecting Piazza Bologna with Conca d'Oro on 13 June 2012; the branch's last stop (and new terminus), Jonio, was opened on 21 April 2015.
Line C
Opened on 9 November 2014, line C currently runs radially from San Giovanni, serving as an interchange station for Line A, to the eastern terminus of Pantano (the former terminus of the Roma–Giardinetti light railway). It is the first Metro line to extend beyond the city boundaries in Rome.
It is planned to extend to the northwest, towards Grottarossa (north of the Vatican), via the city centre; it will also intersect with Line A at Ottaviano (beside the Vatican), with Line B at Colosseo, and with the planned Line D at Piazza Venezia, thus creating a fourth Metro hub in Rome.
The first section of the line, from Centocelle to Pantano, is the furthest from the city centre and includes 15 of the planned 30 stops. A further section of Line C, serving six additional stations, opened on 29 June 2015, as the line's western terminus was moved from Parco di Centocelle to Lodi. On 12 May 2018, the western terminus was moved to San Giovanni (interchange station for line A). After this third phase, the line will be further extended with three stations, Porta Metronia, Colosseo, and Piazza Venezia, located in the city centre.
Progress on the line has been slow with projected completion dates being repeatedly delayed. Rome is one of the oldest cities in the world, and as such, the construction of the Metro system has encountered considerable obstacles owing to frequent archaeological discoveries. While the excavation of the tunnels themselves can be undertaken well below the probable location of most archaeological finds, the excavation of stairwells and ventilation shafts – which must, by necessity, connect with the surface – pose considerable difficulties.
The trains operating on line C are completely automated, and use the same AnsaldoBreda Driverless Metro system also featured on the Copenhagen Metro.
Archeostation
During the excavation for the central route of line C, thanks to the archeological richness of Rome's ground a new type of underground station was born, as in Paris with Louvre–Rivoli station.
San Giovanni
San Giovanni station was the first archeostation to be opened on 12 May 2018. Excavation to a depth of about 20 metres allowed the exploration of about 21 stratifications of history up to the so-called virgin soil, the one in which man's presence is absent.
The exhibition is characterized by being a real tour with libraries for the finds along the route, explanatory panels on the walls and a temporal measurement of the historical phases that follows the path of the passengers from the atrium level to the platforms' level:
Atrium level: from Contemporary Age to Late Antiquity Age
First underground level: from Republican Age to Archaic Age
Platform level: Prehistoric Age
The various archaeological finds and exhibits include small items such as gold jewelry, coins, crockery, shells, large amphorae and elements of ancient columns and also large finds, such as the large pool, the largest reservoir ever found, located inside a farm of the imperial age.
Timeline
Network map
Rolling stock
All Rome metro lines are heavy rapid transit lines, with 6-car trains, approximately 105 m long. Line A of the Rome metro uses exclusively the CAF MA 300 series, line B essentially uses the CAF MB400 series together with other CAF MA300 series trains and the historic MB 100 Ansaldobreda. Line C is the longest driverless metro in Italy and one of the largest in Europe, using Hitachi Rail Italy's driverless technology.
Service
Fares
An urban single journey ticket (integrated ticket), the Biglietto Integrato a Tempo (BIT), costs €1.50, and is valid on the Metro, buses, trams and suburban trains inside the Rome municipality for 100 minutes from first validation. Other tickets are available, including daily (€7.00), 2-day (€12.50), and 3-day (€18.00) passes (Rome 24h/48h/72h), and a weekly pass (€24.00), the Carta Integrata Settimanale (CIS). Monthly passes that are valid during the charged calendar month for unlimited journeys available for the personal usage (€35.00) or impersonal usage (€53.00) and may be used alternatively by different persons. Children under 10 years old travel can travel for free on public transport services when accompanied by a fare-paying adult.
Two proximity cards are available in Rome, ”èRoma” and ”Metrebus Card Red”, which can be charged with season tickets, replacing paper for this type of ticket.
Operating hours
Service starts at about 5:30 am and ends at about 11:30 pm. On Fridays and Saturdays service ends later, at about 1:30 am.
When the Metro is closed, a night bus service operates with lines that follow the same routes and stop at the same stations as the Metro. Line A is served by bus NMA, Line B is served by bus NMB/NMB1, and Line C is served by bus NMC.
Other rail lines
Rome's local transport provider, ATAC, operates the Metro network and the Rome-Giardinetti line. The Roma–Lido, which connects Rome to Ostia, and the Roma–Viterbo line, used to be operated by ATAC until 1 July 2022, when it became part of the Cotral network.
Roma–Lido
Construction of the Roma–Lido line began shortly after the end of World War I and was completed some six years later in 1924. It began operation as a steam locomotion railway but electrification was completed less than a year later.
The line is operated as an integrated part of the Metro, but runs entirely overground. It runs from the Roma Porta San Paolo station beside the Line B Piramide station and runs alongside Line B as far as EUR Magliana. It then continues separately on to the seaside district of Ostia. The line terminates beside the end of Via Cristoforo Colombo.
Roma–Giardinetti
Officially termed a railway, the Roma–Giardinetti line is a narrow gauge tram which connects Laziali (a smaller, local train station some 800 metres east of Termini's main concourse) with Giardinetti, just past the Grande Raccordo Anulare (GRA) – Rome's orbital motorway. The line originally ran to Frosinone some from Rome, but has been gradually reduced in length, when the section from Giardinetti to Pantano, which will become a permanent part of Line C, was taken out of service. Most recently also the part from Centocelle to Giardinetti was reduced.
As regards future projects, the transformation of the line into a modern light rail line is planned, with an extension planned to the Tor Vergata area. The new light rail system will be called "line G" with respect to the metro network and line 11 with respect to the tram network. The work, after a long bureaucratic process, has been commissioned and financed, the line is expected to close during 2026 in order to begin work.
Roma–Civita Castellana–Viterbo
The Roma–Civita Castellana–Viterbo line (also called Roma Nord railway) began life as a narrow-gauge tram running from Piazza della Libertà in Rome to Civita Castellana. However, the next stretch of the line, to Viterbo, was built as a railway and over the years the tram section was converted into a railway as well, a process which concluded with the moving of the Roman terminus from the street-level terminus at Piazza della Libertà across the river to a new underground station in Piazzale Flaminio, beside the subsequently constructed Line A station, after World War II.
The line is operated in two modes: as an urban service from Piazzale Flaminio to Montebello, and as a suburban service from Piazzale Flaminio to Viterbo. The urban service operates with a frequency of about one train every ten minutes, while the suburban service operates considerably less frequently, with less than a third of the trains making the full two-and-a-half–hour journey from Rome to Viterbo.
Future expansions
The Metro system is currently expanding:
An extension of line A from Battistini to Torrevecchia towards the west of the city is planned; a track of 2 kilometres (1.2 mi) with 2 stations.
An extension of line B is also planned. A track of 2.8 kilometres (1.7 mi) with 2 stations should be realized from Rebibbia (Rome Metro) to Torraccia/Casal Monastero, towards the east of Rome. The extension of the line B1 is extending from Jonio to Bufalotta with 3.8 kilometers with 3 stations.
The extension of line C towards the centre of Rome is under construction. From San Giovanni (interchange station with line A) to Grottarossa (Rome Metro), the track is 3.6 kilometres (2.2 mi) long with an intermediate station near San Giovanni Addolorata Hospital, Amba Aradam.
New lines
There are two proposed new lines, with no timeline for their construction:
is currently proposed to run from Ojetti in the north to Agricoltura in the south, passing through the city centre to the west of . Development on the project was suspended from 2012 to 2018.
would take over the existing Rome–Lido railway and the Jonio branch of Line B, with a new railway connection south of /Piramide.
Signalling
Signalling of the Rome Metro guarantees trains' safe and correct movements.
Line A
Line A uses an evolution of the RS4 Codici, a classical block system of the Italian railway. Since its creation the signalling offers to the conductor advice about the speed limit and the freedom of the way.
Line B
Until 1990, Line B used a railway-like signaling system that advised only about the freedom of the way.
Since 1990, the line uses a new signaling system, inspired by the Milan Metro's signaling, that gives information about the speed limit in the section within a range of 0 km/h and 80 km/h. Conductors are informed by classical light semaphores.
Line C
Line C is an automatic line and uses a radio frequency system for communication with the train. It has an electrical block system that permits a max frequency of a train every 90 seconds.
| Technology | Italy | null |
1419898 | https://en.wikipedia.org/wiki/Ribbon%20eel | Ribbon eel | The ribbon eel (Rhinomuraena quaesita), also known as the leaf-nosed moray eel or bernis eel, is a species of moray eel, the only member of the genus Rhinomuraena. The ribbon eel is found in sand burrows and reefs in the Indo-Pacific Ocean. Although generally placed in the moray eel family Muraenidae, it has several distinctive features leading some to place it in its own family, Rhinomuraenidae.
Development
Ribbon eels are a diurnal species (active during the daytime). Characterized by its long, thin body and high dorsal fins. The ribbon eel can easily be recognized by its expanded anterior nostrils and wide-open jaws. These eels can have up to 255 vertebrae in their back bone, making them one of the most narrow and elongated eel species known. The ribbon eel larva is described as a large, greenish leptocephalus. Based on observed color changes, it is generally considered a protandic hermaphrodite (male to female transformation). When necessary, male ribbon eels will develop reproductive organs, lay eggs, and then die. This developmental stage happens within about a month. In their juvenile and subadult stages, ribbon eels appear primarily jet black with a yellow dorsal fin. In adult males, the black is replaced by a vibrant blue and yellow facial appearance. An adult female is entirely yellow or yellow with some blue to the posterior. Color change related to sex change is not known from any other moray eel species. The blue adult males range from in length, while the larger yellow females can reach up to . It is presumed that even with its significant color changes throughout its development, coloring does not play a significant role in mating for a ribbon eel because the eels are colorblind (possessing only one of the two photoreceptor cells required to see colors). In its natural habitat, ribbon eels can live up to twenty years.
Habitat
Ribbon eels prefer more shallow-water areas compared to other moray eels, frequenting a depth range of 1 to 57 meters. This species is widely distributed and are seen by divers in Indonesian waters with their heads and anterior bodies protruding from crevices in sand and rubble habitats, like coral reefs, which they are able to slip through with their slime coat. Typically, ribbon eels can be found in tropical parts of the Indian and Pacific Ocean, ranging from East Africa to southern Japan, Australia, and French Polynesia.
Taxonomy
The ribbon eel is the only one of its genus, Rhinomuraena, in the subspecies Muraeninae of the moray eel family. Its unique morphology, such as its gonad structure and function are not shared with many of the other major genera in the family and provided evidence for separation. Typically, eels of the family urogenital system follow a posterior pattern toward the anus of the eel, but the Rhinomuraena system has its sex cells move anteriorly. This change is presumed to be due to some evolutionary cause, like peripatric speciation. Previously, R. quaesita was used to identify the blue ribbon eels and R. amboinensis was used for black ribbon eels, but these are now recognized as the same species, R. quaesita.
Environmental Impacts
Ribbon eels are labeled "least concerned" on the IUCN Red List of Threatened Species as of 2009. Recent studies on the effects of various environmental factors on marine eels have presented potential threats to the ribbon eel population. Marine eels that prefer to stay in shallower waters, specifically around coral reef habitats, are most at risk of suffering from the impacts of coral bleaching and other marine habitat damage caused by climate change or human activity (i.e. pollution). The shift in habitat structure and diversity could make it harder for ribbon eels to find enough prey to survive and reproduce. Overfishing is another concern as it can lead to a reduction in the number of ribbon eels and their available prey.
Captivity
Most ribbon eels do not live longer than a year in captivity. Ribbon eels have been observed in many cases to stop eating after being captured and put into home aquariums. Higher levels of success have been achieved in public aquaria, where there are a few reported cases of spawning at facilities in Europe and North America. In captivity, the color differences are not related to maturity or gender.
Although captured for the aquarium industry, it remains common and widespread, and is not considered threatened.
| Biology and health sciences | Anguilliformes | Animals |
1420406 | https://en.wikipedia.org/wiki/Sex%20linkage | Sex linkage | Sex linked describes the sex-specific reading patterns of inheritance and presentation when a gene mutation (allele) is present on a sex chromosome (allosome) rather than a non-sex chromosome (autosome). In humans, these are termed X-linked recessive, X-linked dominant and Y-linked. The inheritance and presentation of all three differ depending on the sex of both the parent and the child. This makes them characteristically different from autosomal dominance and recessiveness.
There are many more X-linked conditions than Y-linked conditions, since humans have several times as many genes on the X chromosome than the Y chromosome. Only females are able to be carriers for X-linked conditions; males will always be affected by any X-linked condition, since they have no second X chromosome with a healthy copy of the gene. As such, X-linked recessive conditions affect males much more commonly than females.
In X-linked recessive inheritance, a son born to a carrier mother and an unaffected father has a 50% chance of being affected, while a daughter has a 50% chance of being a carrier, however a fraction of carriers may display a milder (or even full) form of the condition due to a phenomenon known as skewed X-inactivation, in which the normal process of inactivating half of the female body's X chromosomes preferably targets a certain parent's X chromosome (the father's in this case). If the father is affected, the son will not be affected, as he does not inherit the father's X chromosome, but the daughter will always be a carrier (and may occasionally present with symptoms due to aforementioned skewed X-inactivation).
In X-linked dominant inheritance, a son or daughter born to an affected mother and an unaffected father both have a 50% chance of being affected (though a few X-linked dominant conditions are embryonic lethal for the son, making them appear to only occur in females). If the father is affected, the son will always be unaffected, but the daughter will always be affected. A Y-linked condition will only be inherited from father to son and will always affect every generation.
The inheritance patterns are different in animals that use sex-determination systems other than XY. In the ZW sex-determination system used by birds, the mammalian pattern is reversed, since the male is the homogametic sex (ZZ) and the female is heterogametic (ZW).
In classical genetics, a mating experiment called a reciprocal cross is performed to test if an animal's trait is sex-linked.
X-linked dominant inheritance
Each child of a mother affected with an X-linked dominant trait has a 50% chance of inheriting the mutation and thus being affected with the disorder. If only the father is affected, 100% of the daughters will be affected, since they inherit their father's X chromosome, and 0% of the sons will be affected, since they inherit their father's Y chromosome.
There are fewer X-linked dominant conditions than X-linked recessive, because dominance in X-linkage requires the condition to present in females with only a fraction of the reduction in gene expression of autosomal dominance, since roughly half (or as many as 90% in some cases) of a particular parent's X chromosomes are inactivated in females.
Examples
Alport syndrome
Coffin–Lowry syndrome (CLS)
Fragile X syndrome
Idiopathic hypoparathyroidism
Incontinentia pigmenti
Rett syndrome (RS)
Vitamin D resistant rickets (X-linked hypophosphatemia)
X-linked recessive inheritance
Females possessing one X-linked recessive mutation are considered carriers and will generally not manifest clinical symptoms of the disorder, although differences in X chromosome inactivation can lead to varying degrees of clinical expression in carrier females since some cells will express one X allele and some will express the other. All males possessing an X-linked recessive mutation will be affected, since males have only a single X chromosome and therefore have only one copy of X-linked genes. All offspring of a carrier female have a 50% chance of inheriting the mutation if the father does not carry the recessive allele. All female children of an affected father will be carriers (assuming the mother is not affected or a carrier), as daughters possess their father's X chromosome. If the mother is not a carrier, no male children of an affected father will be affected, as males only inherit their father's Y chromosome.
The incidence of X-linked recessive conditions in females is the square of that in males: for example, if 1 in 20 males in a human population are red–green color blind, then 1 in 400 females in the population are expected to be color-blind (1/20)*(1/20).
Examples
Aarskog–Scott syndrome
Adrenoleukodystrophy (ALD)
Bruton's agammaglobulinemia
Color blindness
Complete androgen insensitivity syndrome
Congenital aqueductal stenosis (hydrocephalus)
Duchenne muscular dystrophy
Fabry disease
Glucose-6-phosphate dehydrogenase deficiency
Haemophilia A and B
Hunter syndrome
Inherited nephrogenic diabetes insipidus
Menkes disease (kinky hair syndrome)
Ornithine carbamoyltransferase deficiency
Wiskott–Aldrich syndrome
Y-linked
Various failures in the SRY genes
Sex-linked traits in other animals
White eyes in Drosophila melanogaster flies was one of the earliest sex-linked genes discovered.
Fur color in domestic cats: the gene that causes orange pigment is on the X chromosome; thus a Calico or tortoiseshell cat, with both black (or gray) and orange pigment, is nearly always female.
The first sex-linked gene ever discovered was the "lacticolor" X-linked recessive gene in the moth Abraxas grossulariata by Leonard Doncaster.
Related terms
It is important to distinguish between sex-linked characters, which are controlled by genes on sex chromosomes, and two other categories.
Sex-influenced traits
Sex-influenced or sex-conditioned traits are phenotypes affected by whether they appear in a male or female body. Even in a homozygous dominant or recessive female the condition may not be expressed fully. Example: baldness in humans.
Sex-limited traits
These are characters only expressed in one sex. They may be caused by genes on either autosomal or sex chromosomes. Examples: female sterility in Drosophila; and many polymorphic characters in insects, especially in relation to mimicry. Closely linked genes on autosomes called "supergenes" are often responsible for the latter.
| Biology and health sciences | Genetics | Biology |
1421035 | https://en.wikipedia.org/wiki/Giant%20otter%20shrew | Giant otter shrew | The giant otter shrew (Potamogale velox) is a semiaquatic, carnivorous afrotherian mammal. It is found in the main rainforest block of central Africa from Nigeria to Zambia, with a few isolated populations in Kenya and Uganda. It lives in streams, wetlands and slow flowing larger rivers. It is the only species in the genus Potamogale. Otter shrews are most closely related to the tenrecs of Madagascar.
They are nocturnal carnivores that feed on aquatic animals. Despite its name, the giant otter shrew is neither a true shrew (Soricidae) or otter (Lutrinae). The common name refers to their resemblance to otters with their flat face, stiff whiskers, and muscular tails, and to their overall superficial similarity to true shrews.
Description
The giant otter shrew is a mammal superficially similar to an otter in appearance. It is characterized by a long, flat tail, which it uses for swimming by sideways undulation like a fish. It has a muzzle covered with bristles, and flat shielded nostrils. It has dense, soft hair, silky on the tail.
It has small eyes and external ears. Its fur consists of a dense undercoat and coarse guard hairs. It possesses counter-shading with dark brown on its back and whitish or yellowish under parts. The tail is covered with a short, silky coat of fur and is compressed laterally which allow it to swim by horizontal undulations as in fishes and crocodiles. Its legs are short and lack webbing so they are not used for swimming. The hind feet have a flap of skin along the inside that allows them to be held snugly against the body when swimming. There are also two syndactylous (2nd and 3rd toes are fused) toes on the hind feet, used for grooming. On land P. velox is plantigrade. Females have two mammae on the lower abdomen for feeding young.
The mass ranges from to . Head and body length is to , and reaches to with tail.
Geographic range
Giant otter shrews are native to central Africa, from the southern regions of Nigeria (central Rainforest Zone), and then eastward through Equatorial Guinea, Gabon, and the Central African Republic, Chad, the Republic of Congo, the Democratic Republic of the Congo, South Sudan to the northern regions of Angola and Zambia. There is a small population that lives between Uganda and Kenya and the preserved rainforest of Kakamega, Kenya.
Habitat
This species prefers fresh water aquatic microhabitats in the rainforest. Preferred environments include fast flowing rivers, streams, swamps, coastal rivers, and during rainy season some may retreat to small forest pools (altitude range from 0–1,800m). River banks provide good habitats for breeding and nesting. These animals make burrows with an entrance below water level (like otters) and during the day find shelter there and then become active in the afternoon.
Behavior
The giant otter shrew builds burrows among riverbank crevices. It chooses dry leaves with which to line its nest. This is also where breeding takes place. The burrows are frequently changed. When foraging, otter shrews take frequent grooming breaks. When traveling upstream the otter shrew travels on the bank and then swims downstream. The night foraging routine is regular and predictable, and covers up to 800 meters a night. P. velox regularly visits discrete piles of feces that are sheltered and probably used to mark boundaries of territory.
Giant otter shrews are solitary with one shrew occupying between 500 and 1,000 m of stream.
Food habits
P. velox is a nocturnal predator, hunting primarily by touch and scent in and around calm pools. Each dive lasts only seconds. P. velox searches both within the pool and along the bank for prey using the sensitive vibrissae and odor and apparently not eyesight. It prefers areas that have cover to retreat to when it feels threatened. P. velox attacks prey using sharp bites, sometimes pinning the prey with its fore feet, and flipping crabs over to attack their weaker ventral surface. They usually avoid crabs larger than 7 cm across. The prey preference varies among individuals; some prefer crabs; others, frogs or fish. Frogs are eaten headfirst and fish are pulled apart into manageable bits. Prey is consumed on the bank. P. velox also eats insects, mollusks, and freshwater prawns.
In captivity it eats 15–20 crabs per night.
Lifespan
Giant otter shrews fare very poorly in captivity. Captive specimens have been recorded to deteriorate in health very quickly, living only 1–14 days.
Reproduction
Giant otter shrews breed during the wet/rainy season. They give birth to one or two young per litter, once or twice a year. Males move long distances via water in search of mates and it is thought that males rut (or fight) during the wet season.
Conservation status
Currently this species is listed as being of least concern by the IUCN because its declining rate is not significant enough to move to the next category. However it is on the decline. One of the major threats to this species is the soil erosion caused by deforestation especially in Cameroon. While they can tolerate seasonally cloudy streams, streams muddied from erosion and deforestation are little used. Some drown in fishing nets or fish traps, and members of this species have not survived well in captivity. There is ongoing research about the effects of human activity on them. It is also hunted extensively for its skin.
| Biology and health sciences | Other afrotheres | Animals |
27697009 | https://en.wikipedia.org/wiki/API | API | An application programming interface (API) is a connection between computers or between computer programs. It is a type of software interface, offering a service to other pieces of software. A document or standard that describes how to build such a connection or interface is called an API specification. A computer system that meets this standard is said to implement or expose an API. The term API may refer either to the specification or to the implementation.
In contrast to a user interface, which connects a computer to a person, an application programming interface connects computers or pieces of software to each other. It is not intended to be used directly by a person (the end user) other than a computer programmer who is incorporating it into software. An API is often made up of different parts which act as tools or services that are available to the programmer. A program or a programmer that uses one of these parts is said to call that portion of the API. The calls that make up the API are also known as subroutines, methods, requests, or endpoints. An API specification defines these calls, meaning that it explains how to use or implement them.
One purpose of APIs is to hide the internal details of how a system works, exposing only those parts a programmer will find useful and keeping them consistent even if the internal details later change. An API may be custom-built for a particular pair of systems, or it may be a shared standard allowing interoperability among many systems.
The term API is often used to refer to web APIs, which allow communication between computers that are joined by the internet. There are also APIs for programming languages, software libraries, computer operating systems, and computer hardware. APIs originated in the 1940s, though the term did not emerge until the 1960s and 70s.
Purpose
An API opens a software system to interactions from the outside. It allows two software systems to communicate across a boundary an interface using mutually agreed-upon signals. In other words, an API connects software entities together. Unlike a user interface, an API is typically not visible to users. It is an "under the hood" portion of a software system, used for machine-to-machine communication.
A well-designed API exposes only objects or actions needed by software or software developers. It hides details that have no use. This abstraction simplifies programming.
Building software using APIs has been compared to using building-block toys, such as Lego bricks. Software services or software libraries are analogous to the bricks; they may be joined together via their APIs, composing a new software product. The process of joining is called integration.
As an example, consider a weather sensor that offers an API. When a certain message is transmitted to the sensor, it will detect the current weather conditions and reply with a weather report. The message that activates the sensor is an API call, and the weather report is an API response. A weather forecasting app might integrate with a number of weather sensor APIs, gathering weather data from throughout a geographical area.
An API is often compared to a contract. It represents an agreement between parties: a service provider who offers the API and the software developers who rely upon it. If the API remains stable, or if it changes only in predictable ways, developers' confidence in the API will increase. This may increase their use of the API.
History of the term
The term API initially described an interface only for end-user-facing programs, known as application programs. This origin is still reflected in the name "application programming interface." Today, the term is broader, including also utility software and even hardware interfaces.
The idea of the API is much older than the term itself. British computer scientists Maurice Wilkes and David Wheeler worked on a modular software library in the 1940s for EDSAC, an early computer. The subroutines in this library were stored on punched paper tape organized in a filing cabinet. This cabinet also contained what Wilkes and Wheeler called a "library catalog" of notes about each subroutine and how to incorporate it into a program. Today, such a catalog would be called an API (or an API specification or API documentation) because it instructs a programmer on how to use (or "call") each subroutine that the programmer needs.
Wilkes and Wheeler's book The Preparation of Programs for an Electronic Digital Computer contains the first published API specification. Joshua Bloch considers that Wilkes and Wheeler "latently invented" the API, because it is more of a concept that is discovered than invented.
The term "application program interface" (without an -ing suffix) is first recorded in a paper called Data structures and techniques for remote computer graphics presented at an AFIPS conference in 1968. The authors of this paper use the term to describe the interaction of an application—a graphics program in this case—with the rest of the computer system. A consistent application interface (consisting of Fortran subroutine calls) was intended to free the programmer from dealing with idiosyncrasies of the graphics display device, and to provide hardware independence if the computer or the display were replaced.
The term was introduced to the field of databases by C. J. Date in a 1974 paper called The Relational and Network Approaches: Comparison of the Application Programming Interface. An API became a part of the ANSI/SPARC framework for database management systems. This framework treated the application programming interface separately from other interfaces, such as the query interface. Database professionals in the 1970s observed these different interfaces could be combined; a sufficiently rich application interface could support the other interfaces as well.
This observation led to APIs that supported all types of programming, not just application programming. By 1990, the API was defined simply as "a set of services available to a programmer for performing certain tasks" by technologist Carl Malamud.
The idea of the API was expanded again with the dawn of remote procedure calls and web APIs. As computer networks became common in the 1970s and 80s, programmers wanted to call libraries located not only on their local computers, but on computers located elsewhere. These remote procedure calls were well supported by the Java language in particular. In the 1990s, with the spread of the internet, standards like CORBA, COM, and DCOM competed to become the most common way to expose API services.
Roy Fielding's dissertation Architectural Styles and the Design of Network-based Software Architectures at UC Irvine in 2000 outlined Representational state transfer (REST) and described the idea of a "network-based Application Programming Interface" that Fielding contrasted with traditional "library-based" APIs. XML and JSON web APIs saw widespread commercial adoption beginning in 2000 and continuing as of 2021. The web API is now the most common meaning of the term API.
The Semantic Web proposed by Tim Berners-Lee in 2001 included "semantic APIs" that recast the API as an open, distributed data interface rather than a software behavior interface. Proprietary interfaces and agents became more widespread than open ones, but the idea of the API as a data interface took hold. Because web APIs are widely used to exchange data of all kinds online, API has become a broad term describing much of the communication on the internet. When used in this way, the term API has overlap in meaning with the term communication protocol.
Types
Libraries and frameworks
The interface to a software library is one type of API. The API describes and prescribes the "expected behavior" (a specification) while the library is an "actual implementation" of this set of rules.
A single API can have multiple implementations (or none, being abstract) in the form of different libraries that share the same programming interface.
The separation of the API from its implementation can allow programs written in one language to use a library written in another. For example, because Scala and Java compile to compatible bytecode, Scala developers can take advantage of any Java API.
API use can vary depending on the type of programming language involved.
An API for a procedural language such as Lua could consist primarily of basic routines to execute code, manipulate data or handle errors while an API for an object-oriented language, such as Java, would provide a specification of classes and its class methods. Meanwhile, several studies show that most applications that use an API tend to use a small part of the API.
Language bindings are also APIs. By mapping the features and capabilities of one language to an interface implemented in another language, a language binding allows a library or service written in one language to be used when developing in another language. Tools such as SWIG and F2PY, a Fortran-to-Python interface generator, facilitate the creation of such interfaces.
An API can also be related to a software framework: a framework can be based on several libraries implementing several APIs, but unlike the normal use of an API, the access to the behavior built into the framework is mediated by extending its content with new classes plugged into the framework itself.
Moreover, the overall program flow of control can be out of the control of the caller and in the framework's hands by inversion of control or a similar mechanism.
Operating systems
An API can specify the interface between an application and the operating system. POSIX, for example, specifies a set of common APIs that aim to enable an application written for a POSIX conformant operating system to be compiled for another POSIX conformant operating system.
Linux and Berkeley Software Distribution are examples of operating systems that implement the POSIX APIs.
Microsoft has shown a strong commitment to a backward-compatible API, particularly within its Windows API (Win32) library, so older applications may run on newer versions of Windows using an executable-specific setting called "Compatibility Mode".
An API differs from an application binary interface (ABI) in that an API is source code based while an ABI is binary based. For instance, POSIX provides APIs while the Linux Standard Base provides an ABI.
Remote APIs
Remote APIs allow developers to manipulate remote resources through protocols, specific standards for communication that allow different technologies to work together, regardless of language or platform.
For example, the Java Database Connectivity API allows developers to query many different types of databases with the same set of functions, while the Java remote method invocation API uses the Java Remote Method Protocol to allow invocation of functions that operate remotely, but appear local to the developer.
Therefore, remote APIs are useful in maintaining the object abstraction in object-oriented programming; a method call, executed locally on a proxy object, invokes the corresponding method on the remote object, using the remoting protocol, and acquires the result to be used locally as a return value.
A modification of the proxy object will also result in a corresponding modification of the remote object.
Web APIs
Web APIs are the defined interfaces through which interactions happen between an enterprise and applications that use its assets, which also is a Service Level Agreement (SLA) to specify the functional provider and expose the service path or URL for its API users. An API approach is an architectural approach that revolves around providing a program interface to a set of services to different applications serving different types of consumers.
When used in the context of web development, an API is typically defined as a set of specifications, such as Hypertext Transfer Protocol (HTTP) request messages, along with a definition of the structure of response messages, usually in an Extensible Markup Language (XML) or JavaScript Object Notation (JSON) format. An example might be a shipping company API that can be added to an eCommerce-focused website to facilitate ordering shipping services and automatically include current shipping rates, without the site developer having to enter the shipper's rate table into a web database. While "web API" historically has been virtually synonymous with web service, the recent trend (so-called Web 2.0) has been moving away from Simple Object Access Protocol (SOAP) based web services and service-oriented architecture (SOA) towards more direct representational state transfer (REST) style web resources and resource-oriented architecture (ROA). Part of this trend is related to the Semantic Web movement toward Resource Description Framework (RDF), a concept to promote web-based ontology engineering technologies. Web APIs allow the combination of multiple APIs into new applications known as mashups.
In the social media space, web APIs have allowed web communities to facilitate sharing content and data between communities and applications. In this way, content that is created in one place dynamically can be posted and updated to multiple locations on the web. For example, Twitter's REST API allows developers to access core Twitter data and the Search API provides methods for developers to interact with Twitter Search and trends data.
Design
The design of an API has significant impact on its usage. The principle of information hiding describes the role of programming interfaces as enabling modular programming by hiding the implementation details of the modules so that users of modules need not understand the complexities inside the modules. Thus, the design of an API attempts to provide only the tools a user would expect. The design of programming interfaces represents an important part of software architecture, the organization of a complex piece of software.
Release policies
APIs are one of the more common ways technology companies integrate. Those that provide and use APIs are considered as being members of a business ecosystem.
The main policies for releasing an API are:
Private: The API is for internal company use only.
Partner: Only specific business partners can use the API. For example, vehicle for hire companies such as Uber and Lyft allow approved third-party developers to directly order rides from within their apps. This allows the companies to exercise quality control by curating which apps have access to the API, and provides them with an additional revenue stream.
Public: The API is available for use by the public. For example, Microsoft makes the Windows API public, and Apple releases its API Cocoa, so that software can be written for their platforms. Not all public APIs are generally accessible by everybody. For example, Internet service providers like Cloudflare or Voxility, use RESTful APIs to allow customers and resellers access to their infrastructure information, DDoS stats, network performance or dashboard controls. Access to such APIs is granted either by “API tokens”, or customer status validations.
Public API implications
An important factor when an API becomes public is its "interface stability". Changes to the API—for example adding new parameters to a function call—could break compatibility with the clients that depend on that API.
When parts of a publicly presented API are subject to change and thus not stable, such parts of a particular API should be documented explicitly as "unstable". For example, in the Google Guava library, the parts that are considered unstable, and that might change soon, are marked with the Java annotation @Beta.
A public API can sometimes declare parts of itself as deprecated or rescinded. This usually means that part of the API should be considered a candidate for being removed, or modified in a backward incompatible way. Therefore, these changes allow developers to transition away from parts of the API that will be removed or not supported in the future.
Client code may contain innovative or opportunistic usages that were not intended by the API designers. In other words, for a library with a significant user base, when an element becomes part of the public API, it may be used in diverse ways.
On February 19, 2020, Akamai published their annual “State of the Internet” report, showcasing the growing trend of cybercriminals targeting public API platforms at financial services worldwide. From December 2017 through November 2019, Akamai witnessed 85.42 billion credential violation attacks. About 20%, or 16.55 billion, were against hostnames defined as API endpoints. Of these, 473.5 million have targeted financial services sector organizations.
Documentation
API documentation describes what services an API offers and how to use those services, aiming to cover everything a client would need to know for practical purposes.
Documentation is crucial for the development and maintenance of applications using the API.
API documentation is traditionally found in documentation files but can also be found in social media such as blogs, forums, and Q&A websites.
Traditional documentation files are often presented via a documentation system, such as Javadoc or Pydoc, that has a consistent appearance and structure.
However, the types of content included in the documentation differs from API to API.
In the interest of clarity, API documentation may include a description of classes and methods in the API as well as "typical usage scenarios, code snippets, design rationales, performance discussions, and contracts", but implementation details of the API services themselves are usually omitted. It can take a number of forms, including instructional documents, tutorials, and reference works. It'll also include a variety of information types, including guides and functionalities.
Restrictions and limitations on how the API can be used are also covered by the documentation. For instance, documentation for an API function could note that its parameters cannot be null, that the function itself is not thread safe. Because API documentation tends to be comprehensive, it is a challenge for writers to keep the documentation updated and for users to read it carefully, potentially yielding bugs.
API documentation can be enriched with metadata information like Java annotations. This metadata can be used by the compiler, tools, and by the run-time environment to implement custom behaviors or custom handling.
It is possible to generate API documentation in a data-driven manner. By observing many programs that use a given API, it is possible to infer the typical usages, as well the required contracts and directives. Then, templates can be used to generate natural language from the mined data.
Dispute over copyright protection for APIs
In 2010, Oracle Corporation sued Google for having distributed a new implementation of Java embedded in the Android operating system. Google had not acquired any permission to reproduce the Java API, although permission had been given to the similar OpenJDK project. Judge William Alsup ruled in the Oracle v. Google case that APIs cannot be copyrighted in the U.S. and that a victory for Oracle would have widely expanded copyright protection to a "functional set of symbols" and allowed the copyrighting of simple software commands:
Alsup's ruling was overturned in 2014 on appeal to the Court of Appeals for the Federal Circuit, though the question of whether such use of APIs constitutes fair use was left unresolved.
In 2016, following a two-week trial, a jury determined that Google's reimplementation of the Java API constituted fair use, but Oracle vowed to appeal the decision. Oracle won on its appeal, with the Court of Appeals for the Federal Circuit ruling that Google's use of the APIs did not qualify for fair use. In 2019, Google appealed to the Supreme Court of the United States over both the copyrightability and fair use rulings, and the Supreme Court granted review. Due to the COVID-19 pandemic, the oral hearings in the case were delayed until October 2020.
The case was decided by the Supreme Court in Google's favor.
Examples
ASPI for SCSI device interfacing
Cocoa and Carbon for the Macintosh
DirectX for Microsoft Windows
EHLLAPI
Java APIs
ODBC for Microsoft Windows
OpenAL cross-platform sound API
OpenCL cross-platform API for general-purpose computing for CPUs & GPUs
OpenGL cross-platform graphics API
OpenMP API that supports multi-platform shared memory multiprocessing programming in C, C++, and Fortran on many architectures, including Unix and Microsoft Windows platforms.
Server Application Programming Interface (SAPI)
Simple DirectMedia Layer (SDL)
| Technology | Software development: General | null |
28981081 | https://en.wikipedia.org/wiki/LibreOffice | LibreOffice | LibreOffice () is a free and open-source office productivity software suite, a project of The Document Foundation (TDF). It was forked in 2010 from OpenOffice.org, an open-sourced version of the earlier StarOffice. It consists of programs for word processing; creating and editing spreadsheets, slideshows, diagrams, and drawings; working with databases; and composing mathematical formulae. It is available in 120 languages. Although TDF does not provide support for LibreOffice, enterprise-focused editions are available from companies in the ecosystem.
LibreOffice uses the OpenDocument standard as its native file format but supports formats of most other major office suites, including Microsoft Office, through a variety of import and export filters. It is available for a variety of computing platforms, with official support for Microsoft Windows, macOS, and Linux, and community builds for many other platforms. Ecosystem partner Collabora uses LibreOffice upstream code and provides apps for Android, iOS, iPadOS, and ChromeOS. LibreOffice is the default office suite of the most popular Linux distributions.
LibreOffice Online is an online office suite that includes the applications Writer, Calc, and Impress, and provides an upstream for projects such as commercial Collabora Online. It is the most actively developed free and open-source office suite, with approximately 50 times the development activity of Apache OpenOffice, the other major descendant of OpenOffice.org, in 2015.
The project was announced, and a beta was released on September 28, 2010. LibreOffice was downloaded about 7.5 million times between January 2011 (the first stable release) and October 2011. The project claimed 120 million unique downloading addresses from May 2011 to May 2015 (excluding Linux distributions), with 55 million of those from May 2014 to May 2015. The Document Foundation estimates that there are 200 million active LibreOffice users worldwide, about 25% of whom are students and 10% are Linux users.
Features
Included applications in LibreOffice
Operating systems and processor architectures
LibreOffice is cross platform software. The Document Foundation developers target Microsoft Windows (IA-32 and x86-64), Linux (IA-32, x86-64, and ARM) and macOS (x86-64 and ARM). There are community ports for FreeBSD, NetBSD, OpenBSD and Mac OS X 10.5 PowerPC receive support from contributors to those projects, respectively. LibreOffice is also installable on OpenIndiana via SFE.
Historically, predecessors of LibreOffice, dating back to StarOffice 3, have run on Solaris with SPARC CPUs that Sun Microsystems (and later Oracle) made. Unofficial ports of LibreOffice, whose versions are now obsolete, have supported SPARC. Current unofficial ports of LibreOffice 5.2.5 run only on Intel-compatible hardware, up to Solaris 11.
In 2011, developers announced plans to port LibreOffice both to Android and to iOS. A beta version of a document viewer for Android 4.0 or newer was released in January 2015; in May 2015, LibreOffice Viewer for Android was released with basic editing capabilities. In February 2020, Collabora released its first officially supported version of LibreOffice (branded as Collabora Office) for Android and iOS.
In July 2020, Collabora shipped an app, branded Collabora Office, for ChromeOS, as used on the popular Chromebook line of notebook computers as well as other form factors of computers. The LibreOffice Impress Remote application for various mobile operating systems allows for remote control of LibreOffice Impress presentations. In June 2023, Red Hat announced that it will no longer support LibreOffice on future editions of Red Hat Enterprise Linux in order to focus on Wayland support and other priorities towards workstation users. LibreOffice will still be available via distribution-neutral Flatpak. Starting with LibreOffice 7.6 on Fedora 39, packaging and maintenance of LibreOffice on Fedora Linux will be managed by the Fedora LibreOffice Special Interest Group instead of Red Hat.
Table of cross platform support
LibreOffice Online
LibreOffice Online is the online office suite edition of LibreOffice. It allows for the use of LibreOffice through a web browser by using the canvas element of HTML5. Development was announced at the first LibreOffice Conference in October 2011, and is ongoing. The Document Foundation, IceWarp, and Collabora announced a collaboration to work on its implementation. A version of the software was shown in a September 2015 conference, and the UK Crown Commercial Service announced an interest in using the software. On 15 December 2015, Collabora, in partnership with ownCloud, released a technical preview of LibreOffice Online branded as Collabora Online Development Edition (CODE).
In July 2016, the enterprise version Collabora Online 1.0 was released. That same month, Nextcloud and Collabora partnered to bring CODE to Nextcloud users. By October 2016, Collabora had released nine updates to CODE. The first source code release of LibreOffice Online was done with LibreOffice version 5.3 in February 2017. In June 2019, CIB software GmbH officially announced its contributions to LibreOffice Online and "LibreOffice Online powered by CIB". In October 2020, Collabora announced the move of its work on Collabora Online from The Document Foundation infrastructure to GitHub.
Comparison with OpenOffice
A detailed 60-page report in June 2015 compared the progress of the LibreOffice project with the related project Apache OpenOffice. It showed that "OpenOffice received about 10% of the improvements LibreOffice did in the period of time studied."
Supported file formats
As its native file format to save documents for all of its applications, LibreOffice uses the Open Document Format for Office Applications (ODF), or OpenDocument, an international standard developed jointly by the International Organization for Standardization (ISO) and the International Electrotechnical Commission (IEC). LibreOffice also supports the file formats of most other major office suites, including Microsoft Office, through a variety of import and export filters.
Miscellaneous features
LibreOffice can use the GStreamer multimedia framework in Linux to render multimedia content such as videos in Impress and other programs. Visually, LibreOffice used the large "Tango style" icons that are used for the application shortcuts, quick launch icons, icons for associated files and for the icons found on the toolbar of the LibreOffice programs in the past, and used on the toolbars and menus by default. They were later replaced by multiple icon themes to adapt the look and feel of specific desktop environment, such as Colibre for Windows, and Elementary for GNOME.
LibreOffice also ships with a modified theme which looks native on GTK-based Linux distributions. It also renders fonts via Cairo on Linux distributions; this means that text in LibreOffice is rendered the same as the rest of the Linux desktop. With version 6.2, LibreOffice includes a ribbon-style GUI, called Notebookbar, including three different views. This feature has formerly been included as an experimental feature in LibreOffice 6 (experimental features must be enabled from LibreOffice settings to make the option available in the View menu).
LibreOffice has a feature similar to WordArt called Fontwork. LibreOffice uses HarfBuzz for complex text layout, it was first introduced in 4.1 for Linux and 5.3 for Windows and macOS. Fonts with OpenType, Apple Advanced Typography or SIL Graphite features can be switched by either a syntax in the Font Name input box, or the Font Features dialog from the Character dialog. LibreOffice supports a "hybrid PDF" format, a file in Portable Document Format (PDF) which can be read by any program supporting PDF, but also contains the source document in ODF format, editable in LibreOffice by dragging and dropping.
Licensing
The LibreOffice project uses a dual LGPLv3 (or later) / MPL 2.0 license for new contributions to allow the license to be upgraded. Since the core of the OpenOffice.org codebase was donated to the Apache Software Foundation, there is an ongoing effort to get all the code rebased to ease future license updates.
Scripting and extensions
LibreOffice supports third-party extensions. , the LibreOffice Extension Repository lists more than 320 extensions. Another list is maintained by the Apache Software Foundation and another one by the Free Software Foundation. Extensions and scripts for LibreOffice can be written in C++, Java, CLI, Python, and LibreOffice Basic. Interpreters for the latter two are bundled with most LibreOffice installers, so no additional installation is needed. The application programming interface for LibreOffice is called "UNO" and is extensively documented.
LibreOffice Basic
LibreOffice Basic is a programming language similar to Microsoft Visual Basic for Applications (VBA) but based on StarOffice Basic. It is available in Writer, Calc and Base. It is used to write small programs known as "macros", with each macro performing a different task, such as counting the words in a paragraph.
History
ooo-build, Go-oo and Oracle
Members of the OpenOffice.org community who were not Sun Microsystems employees had wanted a more egalitarian form for the OpenOffice.org project for many years; Sun had stated in the original OpenOffice.org announcement in 2000 that the project would eventually be run by a neutral foundation, and put forward a more detailed proposal in 2001. Ximian and then Novell had maintained the ooo-build patch set, a project led by Michael Meeks, to make the build easier on Linux and because it was difficult to get contributions accepted by Sun, even from corporate partners. It tracked the main line of development and was not intended to constitute a fork. It was also the standard build mechanism for OpenOffice.org in most Linux distributions, and was contributed to by said distributions.
In 2007, ooo-build was made available by Novell as a software package called Go-oo (ooo-build had used the go-oo.org domain name as early as 2005), which included many features not included in upstream OpenOffice.org. Go-oo also encouraged outside contributions, with rules similar to those later adopted for LibreOffice. Sun's contributions to OpenOffice.org had been declining for some time. They remained reluctant to accept contributions and contributors were upset at Sun releasing OpenOffice.org code to IBM for IBM Lotus Symphony under a proprietary contract, rather than under an open source licence.
Sun was purchased by Oracle Corporation in early 2010. OpenOffice.org community members were concerned by Oracle's behaviour towards open source software, specifically the Java lawsuit against Google and Oracle's withdrawal of developers, and lack of activity on or visible commitment to OpenOffice.org, as had been noted by industry observers; as Meeks put it in early September 2010, "The news from the Oracle OpenOffice conference was that there was no news." Discussion of a fork started soon after.
The Document Foundation and LibreOffice
On 28 September 2010, The Document Foundation was announced as the host of LibreOffice, a new derivative of OpenOffice.org. The Document Foundation's initial announcement stated their concerns that Oracle would either discontinue OpenOffice.org, or place restrictions on it as an open source project, as it had on Sun's OpenSolaris. LibreOffice 3.3 beta used the ooo-build build infrastructure and the OpenOffice.org 3.3 beta code from Oracle, then adding selected patches from Go-oo. Go-oo was discontinued in favour of LibreOffice. Since the office suite that was branded "OpenOffice.org" in most Linux distributions was in fact Go-oo, most moved immediately to LibreOffice. Oracle was invited to become a member of The Document Foundation; however, Oracle demanded that all members of the OpenOffice.org Community Council involved with The Document Foundation step down from the OOo Community Council, claiming a conflict of interest.
Naming
The name "LibreOffice" was picked after research in trademark databases and social media and checks to ensure it could be used for URLs in various countries. Oracle rejected requests to donate the OpenOffice.org brand to the project. LibreOffice was initially named BrOffice in Brazil. OpenOffice.org had been distributed as BrOffice.org by the BrOffice Centre of Excellence for Free Software because of a trademark issue.
End of OpenOffice.org and beginning of Apache OpenOffice
Oracle announced in April 2011 that it was ending its development of OpenOffice.org and would lay off the majority of its paid developers. In June 2011, Oracle announced that it would donate the OpenOffice.org code and trademark to the Apache Software Foundation, where the project was accepted for a project incubation process within the foundation, thus becoming Apache OpenOffice. In an interview with LWN in May 2011, Ubuntu founder Mark Shuttleworth blamed The Document Foundation for destroying OpenOffice.org because it did not license its code under Oracle's Contributor License Agreement. In opposition to Shuttleworth's view, the former Sun executive Simon Phipps argued in the interview for the same online magazine, that the lay-off was an inevitable business decision by Oracle, not impacted by existence of LibreOffice. In March 2015, an LWN.net comparison of LibreOffice with its cousin project Apache OpenOffice concluded that "LibreOffice has won the battle for developer participation".
Release history
Mascot competition
In late 2017 The Document Foundation held a competition for the new mascot of LibreOffice. The mascot was to be used primarily by the community, and was not intended to supersede existing logos for the project. Over 300 concepts were submitted before the first evaluation phase. The mascot contest was cancelled soon after new submissions stopped being accepted. The Document Foundation cited their lack of clear rules and arguments among community members as their reasoning for cancelling the contest.
Versions
Since March 2014 and version 4.2.2, two different major "released" versions of LibreOffice are available at any time in addition to development versions (numbered release candidates and dated nightly builds). The versions are designated to signal their appropriateness for differing user requirements. Releases are designated by three numbers separated by dots. The first two numbers represent the major version (branch) number, and the final number indicates the bugfix releases made in that series. LibreOffice designates the two release versions as:
"Fresh" – the most recent major version (branch), which contains the latest enhancements but which may have introduced bugs not present in the "still" release.
"Still" (formerly "Stable") – the prior major version, which, by the time it has become the "still" version, has had around six months of bug fixing. It is recommended for users for whom stability is more important than the latest enhancements.
Since January 2024 and version 24.2.0, LibreOffice use calendar-based release numbering scheme.
Release schedule
LibreOffice uses a time-based release schedule for predictability, rather than a "when it's ready" schedule. New major versions are released around every six months, in January or February and July or August of each year. The initial intention was to release in March and September, to align with the schedule of other free software projects.
Minor bugfix versions of the "fresh" and "still" release branches are released frequently.
Enterprise support
Commercially supported distributions for LibreOffice with service-level agreements are available via partners such as Collabora (marketed as Collabora Office and Collabora Online), CIB (marketed as CIB Office on the Microsoft Store), and Red Hat. The three vendors are major corporate contributors to the LibreOffice project. As of version 7.1, the open source release of LibreOffice is officially branded as "LibreOffice Community", in order to emphasize that the releases are intended primarily for personal individual use, and are "not targeted at enterprises, and not optimized for their support needs". The Document Foundation states that usage of the community versions in such settings "has had a two-fold negative consequence for the project: a poor use of volunteers' time, as they have to spend their time to solve problems for business that provide nothing in return to the community, and a net loss for ecosystem companies."
Users and deployments
The figure shows the worldwide number of LibreOffice users from 2011 to 2018 in millions. | Technology | Office and data management | null |
922939 | https://en.wikipedia.org/wiki/Acorn | Acorn | The acorn is the nut of the oaks and their close relatives (genera Quercus, Notholithocarpus and Lithocarpus, in the family Fagaceae). It usually contains a seedling surrounded by two cotyledons (seedling leaves), enclosed in a tough shell known as the pericarp, and borne in a cup-shaped cupule. Acorns are long and on the fat side. Acorns take between 5 and 24 months (depending on the species) to mature; see the list of Quercus species for details of oak classification, in which acorn morphology and phenology are important factors.
Etymology
The word acorn (earlier akerne, and acharn) is related to the Gothic name akran, which had the sense of "fruit of the unenclosed land". The word was applied to the most important forest produce, that of the oak. Chaucer spoke of "achornes of okes" in the 14th century. By degrees, popular etymology connected the word both with "corn" and "oak-horn", and the spelling changed accordingly. The current spelling (emerged century) derives from association with ac (Old English: "oak") + corn.
Ecology
Acorns play an important role in forest ecology when oaks are plentiful or dominant in the landscape. The volume of the acorn crop may vary widely, creating great abundance or great stress on the many animals dependent on acorns and the predators of those animals. Acorns, along with other nuts, are termed mast.
Wildlife that consume acorns as an important part of their diets include birds, such as jays, pigeons, some ducks, and several species of woodpeckers. Small mammals that feed on acorns include mice, squirrels and several other rodents. One beetle species, Thorectes lusitanicus, also feeds on acorns. Acorns have a large influence on small rodents in their habitats, as large acorn yields help rodent populations to grow.
Large mammals such as pigs, bears, and deer also consume large amounts of acorns; they may constitute up to 25% of the diet of deer in the autumn. In Spain, Portugal and the New Forest region of southern England, pigs are still turned loose in dehesas (large oak groves) in the autumn, to fill and fatten themselves on acorns. Heavy consumption of acorns can, on the other hand, be toxic to other animals that cannot detoxify their tannins, such as horses and cattle, especially if eaten in excess.
The larvae of some moths and weevils also live in young acorns, consuming the kernels as they develop.
Acorns are attractive to animals because they are large and thus efficiently consumed or cached. Acorns are also rich in nutrients. Percentages vary from species to species, but all acorns contain large amounts of protein, carbohydrates and fats, as well as the minerals calcium, phosphorus and potassium, and the vitamin niacin. Total food energy in an acorn also varies by species, but all compare well with other wild foods and with other nuts.
Acorns also contain bitter tannins, the amount varying with the species. Since tannins, which are plant polyphenols, interfere with an animal's ability to metabolize protein, creatures must adapt in different ways to use the nutritional value acorns contain. Animals may preferentially select acorns that contain fewer tannins. When the tannins are metabolized in cattle, the tannic acid produced can cause ulceration and kidney failure.
Animals that cache acorns, such as jays and squirrels, may wait to consume some of these acorns until sufficient groundwater has percolated through them to leach out the tannins. Other animals buffer their acorn diet with other foods. Many insects, birds, and mammals metabolize tannins with fewer ill effects than do humans.
Species of acorn that contain large amounts of tannins are very bitter, astringent, and potentially irritating if eaten raw. This is particularly true of the acorns of American red oaks and English oaks. The acorns of white oaks, being much lower in tannins, are nutty in flavor; this characteristic is enhanced if the acorns are given a light roast before grinding.
Tannins can be removed by soaking chopped acorns in several changes of water, until the water no longer turns brown. Cold water leaching can take several days, but three to four changes of boiling water can leach the tannins in under an hour. Hot water leaching (boiling) cooks the starch of the acorn, which would otherwise act like gluten in flour, helping it bind to itself. For this reason, if the acorns will be used to make flour, then cold water leaching is preferred.
Being rich in fat, acorn flour can spoil or molder easily and must be carefully stored. Acorns are also sometimes prepared as a massage oil.
Acorns of the white oak group, Leucobalanus, typically start rooting as soon as they are in contact with the soil (in the fall), then send up the leaf shoot in the spring.
Dispersal agents
Acorns are too heavy for wind dispersal, so they require other ways to spread. Oaks therefore depend on biological seed dispersal agents to move the acorns beyond the mother tree and into a suitable area for germination (including access to adequate water, sunlight and soil nutrients), ideally a minimum of from the parent tree.
Many animals eat unripe acorns on the tree or ripe acorns from the ground, with no reproductive benefit to the oak, but some animals, such as squirrels and jays serve as seed dispersal agents. Jays and squirrels that scatter-hoard acorns in caches for future use effectively plant acorns in a variety of locations in which it is possible for them to germinate and thrive.
Even though jays and squirrels retain remarkably large mental maps of cache locations and return to consume them, the odd acorn may be lost, or a jay or squirrel may die before consuming all of its stores. A small number of acorns manage to germinate and survive, producing the next generation of oaks.
Scatter-hoarding behavior depends on jays and squirrels associating with plants that provide good packets of food that are nutritionally valuable, but not too big for the dispersal agent to handle. The beak sizes of jays determine how large acorns may get before jays ignore them.
Acorns germinate on different schedules, depending on their place in the oak family. Once acorns sprout, they are less nutritious, as the seed tissue converts to the indigestible lignins that form the root.
Uses
In some cultures, acorns once constituted a dietary staple, though they have largely been replaced by grains and are now typically considered a relatively unimportant food, except in some Native American and Korean communities.
Several cultures have devised traditional acorn-leaching methods, sometimes involving specialized tools, that were traditionally passed on to their children by word of mouth.
Culinary
Acorns served an important role in early human history and were a source of food for many cultures around the world. For instance, the Ancient Greek lower classes and the Japanese (during the Jōmon period) would eat acorns, especially in times of famine. In ancient Iberia they were a staple food, according to Strabo. Despite this history, acorns rarely form a large part of modern diets and are not currently cultivated on scales approaching that of many other nuts. However, if properly prepared (by selecting high-quality specimens and leaching out the bitter tannins in water), acorn meal can be used in some recipes calling for grain flours. In antiquity, Pliny the Elder noted that acorn flour could be used to make bread. Varieties of oak differ in the amount of tannin in their acorns. Varieties preferred by Native Americans, such as Quercus kelloggii (California black oak), may be easier to prepare or more palatable.
In Korea, an edible jelly named dotorimuk is made from acorns, and dotori guksu are Korean noodles made from acorn flour or starch. In the 17th century, a juice extracted from acorns was administered to habitual drunkards to cure them of their condition or else to give them the strength to resist another bout of drinking.
Roasted acorn flour is a main ingredient in sweet cakes special to Kurdish areas of Iran and Iraq.
Acorns have frequently been used as a coffee substitute, particularly when coffee was unavailable or rationed. The Confederates in the American Civil War and Germans during World War I (when it was called Ersatz coffee), which were cut off from coffee supplies by Union and Allied blockades respectively, are particularly notable past instances of this use of acorns.
Use by Native Americans
Acorns are a traditional food of many indigenous peoples of North America, and long served an especially important role for Californian Native Americans, where the ranges of several species of oaks overlap, increasing the reliability of the resource. One ecology researcher of Yurok and Karuk heritage reports that "his traditional acorn preparation is a simple soup, cooked with hot stones directly in a basket," and says he enjoys acorns eaten with "grilled salmon, huckleberries or seaweed."
Unlike many other plant foods, acorns do not need to be eaten or processed right away, but may be stored for a long time, much as squirrels do. In years that oaks produced many acorns, Native Americans sometimes collected enough acorns to store for two years as insurance against poor acorn production years.
After drying in the sun to discourage mould and germination, acorns could be cached in hollow trees or structures on poles to keep them safe from mice and squirrels. Stored acorns could then be used when needed, particularly during the winter when other resources were scarce. Acorns that germinated in the fall were shelled and pulverized before those germinating in spring. Because of their high fat content, stored acorns can become rancid. Moulds may also grow on them.
The lighting of ground fires killed the larvae of acorn moths and acorn weevils by burning them during their dormancy period in the soil. The pests can infest and consume more than 95% of an oak's acorns.
Fires also released the nutrients bound in dead leaves and other plant debris into the soil, thus fertilizing oak trees while clearing the ground to make acorn collection easier. Most North American oaks tolerate light fires, especially when consistent burning has eliminated woody fuel accumulation around their trunks. Consistent burning encouraged oak growth at the expense of other trees less tolerant of fire, thus keeping oaks dominant in the landscapes.
Oaks produce more acorns when they are not too close to other oaks and thus competing with them for sunlight, water and soil nutrients. The fires tended to eliminate the more vulnerable young oaks and leave old oaks which created open oak savannas with trees ideally spaced to maximize acorn production.
In culture
Art
A motif in Roman architecture, also popular in Celtic and Scandinavian art, the acorn symbol is used as an ornament on cutlery, furniture, and jewelry; it also appears on finials at Westminster Abbey.
In the Artemis Fowl book series, "The Ritual" describes the method used by faeries to regenerate their magical powers.
Military symbolism
The acorn was used frequently by both Union and Confederate forces during the American Civil War. Modern US Army Cavalry Scout campaign hats still retain traces of the acorn today.
Contemporary use as symbol
The acorn is the symbol for the National Trails of England and Wales, and is used for the waymarks on these paths. The acorn, specifically that of the white oak, is also present in the symbol for the University of Connecticut.
Acorns are also used as charges in heraldry.
| Biology and health sciences | Nuts | Plants |
923714 | https://en.wikipedia.org/wiki/Antiaromaticity | Antiaromaticity | Antiaromaticity is a chemical property of a cyclic molecule with a π electron system that has higher energy, i.e., it is less stable due to the presence of 4n delocalised (π or lone pair) electrons in it, as opposed to aromaticity. Unlike aromatic compounds, which follow Hückel's rule ([4n+2] π electrons) and are highly stable, antiaromatic compounds are highly unstable and highly reactive. To avoid the instability of antiaromaticity, molecules may change shape, becoming non-planar and therefore breaking some of the π interactions. In contrast to the diamagnetic ring current present in aromatic compounds, antiaromatic compounds have a paramagnetic ring current, which can be observed by NMR spectroscopy.
Examples of antiaromatic compounds are pentalene (A), biphenylene (B), cyclopentadienyl cation (C). The prototypical example of antiaromaticity, cyclobutadiene, is the subject of debate, with some scientists arguing that antiaromaticity is not a major factor contributing to its destabilization.
Cyclooctatetraene is an example of a molecule adopting a non-planar geometry to avoid the destabilization that results from antiaromaticity. If it were planar, it would have a single eight-electron π system around the ring, but it instead adopts a boat-like shape with four individual π bonds. Because antiaromatic compounds are often short-lived and difficult to work with experimentally, antiaromatic destabilization energy is often modeled by simulation rather than by experimentation.
Definition
The term 'antiaromaticity' was first proposed by Ronald Breslow in 1967 as "a situation in which a cyclic delocalisation of electrons is destabilising". The IUPAC criteria for antiaromaticity are as follows:
The molecule must be cyclic.
The molecule must be planar.
The molecule must have a complete conjugated π-electron system within the ring.
The molecule must have 4n π-electrons where n is any integer within the conjugated π-system.
This differs from aromaticity only in the fourth criterion: aromatic molecules have 4n +2 π-electrons in the conjugated π system and therefore follow Hückel’s rule. Non-aromatic molecules are either noncyclic, nonplanar, or do not have a complete conjugated π system within the ring.
Having a planar ring system is essential for maximizing the overlap between the p orbitals which make up the conjugated π system. This explains why being a planar, cyclic molecule is a key characteristic of both aromatic and antiaromatic molecules. However, in reality, it is difficult to determine whether or not a molecule is completely conjugated simply by looking at its structure: sometimes molecules can distort in order to relieve strain and this distortion has the potential to disrupt the conjugation. Thus, additional efforts must be taken in order to determine whether or not a certain molecule is genuinely antiaromatic.
An antiaromatic compound may demonstrate its antiaromaticity both kinetically and thermodynamically. As will be discussed later, antiaromatic compounds experience exceptionally high chemical reactivity. Being highly reactive is not "indicative" of an antiaromatic compound, but merely suggests that the compound could be antiaromatic. An antiaromatic compound may also be recognized thermodynamically by measuring the energy of the cyclic conjugated π electron system. In an antiaromatic compound, the amount of conjugation energy in the molecule will be significantly higher than in an appropriate reference compound.
In reality, it is recommended that one analyze the structure of a potentially antiaromatic compound extensively before declaring that it is indeed antiaromatic. If an experimentally determined structure of the molecule in question does not exist, a computational analysis must be performed. The potential energy of the molecule should be probed for various geometries in order to assess any distortion from a symmetric planar conformation.
This procedure is recommended because there have been multiple instances in the past where molecules which appear to be antiaromatic on paper turn out to be not truly so in actuality. The most famous (and heavily debated) of these molecules is cyclobutadiene, as is discussed later.
Examples of antiaromatic compounds are pentalene (A), biphenylene (B), cyclopentadienyl cation (C). The prototypical example of antiaromaticity, cyclobutadiene, is the subject of debate, with some scientists arguing that antiaromaticity is not a major factor contributing to its destabilization. Cyclooctatetraene appears at first glance to be antiaromatic, but is an excellent example of a molecule adopting a non-planar geometry to avoid the destabilization that results from antiaromaticity. Because antiaromatic compounds are often short-lived and difficult to work with experimentally, antiaromatic destabilization energy is often modeled by simulation rather than by experimentation.
NMR spectroscopy
The paramagnetic ring current resulting from the electron delocalization in antiaromatic compounds can be observed by NMR. This ring current leads to a deshielding (downfield shift) of nuclei inside the ring and a shielding (upfield shift) of nuclei outside the ring. [12]annulene is an antiaromatic hydrocarbon that is large enough to have protons both inside and outside of the ring. The chemical shift for the protons outside its ring is 5.91 ppm and that for the protons inside the ring is 7.86 ppm, compared to the normal range of 4.5-6.5 ppm for nonaromatic alkenes. This effect is of a smaller magnitude than the corresponding shifts in aromatic compounds.
Many aromatic and antiaromatic compounds (benzene and cyclobutadiene) are too small to have protons inside of the ring, where shielding and deshielding effects can be more diagnostically useful in determining if a compound is aromatic, antiaromatic, or nonaromatic. Nucleus Independent Chemical Shift (NICS) analysis is a method of computing the ring shielding (or deshielding) at the center of a ring system to predict aromaticity or antiaromaticity. A negative NICS value is indicative of aromaticity and a positive value is indicative of antiaromaticity.
Examples
While there are multitudes of molecules in existence which would appear to be antiaromatic on paper, the number of molecules that are antiaromatic in actuality is considerably less. This is compounded by the fact that one cannot typically make derivatives of antiaromatic molecules by adding more antiaromatic hydrocarbon rings, etc. because the molecule typically loses either its planar nature or its conjugated system of π-electrons and becomes nonaromatic. In this section, only examples of antiaromatic compounds which are non-disputable are included.
Pentalene is an antiaromatic compound which has been well-studied both experimentally and computationally for decades. It is dicyclic, planar and has eight π-electrons, fulfilling the IUPAC definition of antiaromaticity. Pentalene’s dianionic and dicationic states are aromatic, as they follow Hückel’s 4n +2 π-electron rule.
Hexadehydro-[12]annulene
Like its relative [12]annulene, hexadehydro-[12]annulene is also antiaromatic. Its structure has been studied computationally via ab initio and density functional theory calculations and is confirmed to be antiaromatic.
Cyclobutadiene
Cyclobutadiene is a classic textbook example of an antiaromatic compound. It is conventionally understood to be planar, cyclic, and have 4 π electrons (4n for n=1) in a conjugated system.
However, it has long been questioned if cyclobutadiene is genuinely antiaromatic and recent discoveries have suggested that it may not be. Cyclobutadiene is particularly destabilized and this was originally attributed to antiaromaticity. However, cyclobutadiene adopts more double bond character in two of its parallel bonds than others and the π electrons are not delocalized between the two double-bond-like bonds, giving it a rectangular shape as opposed to a regular square. As such, cyclobutadiene behaves like two discrete alkenes joined by two single bonds, and is therefore non-aromatic rather than antiaromatic.
Despite the lack of this π-antiaromatic destabilization effect, none of its 4n π-electron relatives (cyclooctatetraene, etc.) had even close to as much destabilization, suggesting there was something more going on in the case of cyclobutadiene. It was found that a combination of angle strain, torsional strain, and Pauli repulsion leads to the extreme destabilization experienced in this molecule.
This discovery is awkward in that it contradicts basic teachings of antiaromaticity. At this point of time, it is presumed that cyclobutadiene will continue to be used to introduce the concept of antiaromaticity in textbooks as a matter of convenience, even though classifying it as antiaromatic technically may not be accurate.
Cyclopentadienyl cation
The cyclopentadienyl cation is another textbook example of an antiaromatic compound. It is conventionally understood to be planar, cyclic, and have 4 π electrons (4n for n=1) in a conjugated system.
However, it has long been questioned if the cyclopentadienyl cation is genuinely antiaromatic and recent discoveries have suggested that it may not be. The lowest-energy singlet state is antiaromatic, but the lowest-energy triplet state is aromatic due to Baird's rule, and research in 2007 showed the triplet state to be the ground state.
Cyclooctatetraene
Cyclooctatetraene is another example of a molecule which is not antiaromatic, even though it might initially appear to be so. Cyclooctatetraene assumes a tub (i.e., boat-like) conformation. As it is not planar, even though it has 4n π-electrons, these electrons are not delocalized and conjugated. The molecule is therefore non-aromatic.
Effects on reactivity
Antiaromatic compounds, often being very unstable, can be highly reactive in order to relieve the antiaromatic destabilization. Cyclobutadiene, for example, rapidly dimerizes with no potential energy barrier via a 2 + 2 cycloaddition reaction to form tricyclooctadiene. While the antiaromatic character of cyclobutadiene is the subject of debate, the relief of antiaromaticity is usually invoked as the driving force of this reaction.
Antiaromaticity can also have a significant effect on pKa. The linear compound propene has a pKa of 44, which is relatively acidic for an sp3 carbon center because the resultant allyl anion can be resonance stabilized. The analogous cyclic system appears to have even more resonance stabilized, as the negative charge can be delocalized across three carbons instead of two. However, the cyclopropenyl anion has 4 π electrons in a cyclic system and in fact has a substantially higher pKa than 1-propene because it is antiaromatic and thus destabilized. Because antiaromatic compounds are often short-lived and difficult to work with experimentally, antiaromatic destabilization energy is often modeled by simulation rather than by experimentation.
Some antiaromatic compounds are stable, especially larger cyclic systems (in which the antiaromatic destabilization is not as substantial). For example, the aromatic species 1 can be reduced to 2 with a relatively small penalty for forming an antiaromatic system. The antiaromatic 2 does revert to the aromatic species 1 over time by reacting with oxygen in the air because the aromaticity is preferred.
The loss of antiaromaticity can sometimes be the driving force of a reaction. In the following keto-enol tautomerization, the product enol is more stable than the original ketone even though the ketone contains an aromatic benzene moiety (blue). However, there is also an antiaromatic lactone moiety (green). The relief of antiaromatic destabilization provides a driving force that outweighs even the loss of an aromatic benzene.
| Physical sciences | Aromatic hydrocarbons | Chemistry |
923872 | https://en.wikipedia.org/wiki/Quern-stone | Quern-stone | Quern-stones are stone tools for hand-grinding a wide variety of materials, especially for various types of grains. They are used in pairs. The lower stationary stone of early examples is called a saddle quern, while the upper mobile stone is called a muller, rubber, or handstone. The upper stone was moved in a back-and-forth motion across the saddle quern. Later querns are known as rotary querns. The central hole of a rotary quern is called the eye, and a dish in the upper surface is known as the hopper. A handle slot contained a handle which enabled the rotary quern to be rotated. They were first used in the Neolithic era to grind cereals into flour.
Design of quern-stones
The upper stones were usually concave while the lower ones were convex. Quern-stones are frequently identifiable by their grooved working surfaces which enabled the movement of flour. Sometimes a millrind was present as a piece of wood (or other material), which allowed the cereal etc. to be added but still acted as a centering device. The upper stone sometimes had a cup-shaped area around the hopper hole with a raised edge. Most handstones have a handle hole on the upper surface, but one class of quern-stones have a slot handle which indicates that a piece of wood was placed horizontally and protruded out from the edge so that the operator could turn the stone by standing and using a rod vertically. One class of upper quern-stones has from two to three sockets for the rod used to turn them and this is thought to reflect the need to reduce wear and tear by having alternative points of contact when in active use.
Grain
Quern-stones have been used by numerous civilizations throughout the world to grind materials, the most important of which was usually grain to make flour for bread-making. They were generally replaced by millstones once mechanised forms of milling appeared, particularly the water mill and the windmill, although animals were also used to operate the millstones. However, in many non-Westernised, non-mechanised cultures they are still manufactured and used regularly and have only been replaced in many parts of the world in the last century or so.
The use of grinding stones for vegetal food processing, and possibly the production of flour, was widespread across Europe from at least 30,000 years ago.
In early Maya civilizations the process of nixtamalization was distinctive in that hard, ripe kernels of maize (corn) were boiled in water and lime, thus producing nixtamal which was then made into unleavened dough for flat cakes by grinding with a handstone on a quern (metate).
Quern stones were used in China at least 10,000 years ago to grind wheat into flour. The production of flour by rubbing wheat by hand took several hours. Due to their form, dimensions, and the nature of the treatment of the surfaces, they reproduce precisely the most ancient implements used for grinding cereal grain into flour. Saddle querns were known in China during the Neolithic Age but rotary stone mills did not appear until the Warring States Period. A prehistoric quern dating back to 23,000 BCE was found at the Longwangchan archaeological site, in Hukou, Shaanxi in 2007. The site is located in the heartland of the northern Chinese loess plateau near the Yellow River.
Other materials
As well as grain, ethnographic evidence and Mesopotamian texts show that a wide range of foodstuffs and inorganic materials were processed using stone querns or mortars, including nuts, seeds, fruit, vegetables, herbs, spices, meat, bark, pigments, temper and clay. Moreover, one study analysing quern-stones noted that a number of querns had traces of arsenic and bismuth, unlike their source rocks, and had levels of antimony which were ten times higher than those of the rocks. The authors concluded that this was probably due to the use of these querns in the preparation of medicines, cosmetics, dyes or even in the manufacture of alloys.
Querns were widely used in grinding metals ores after mining extraction. The aim was to liberate fine ore particles which could then be separated by washing for example, prior to smelting. They were thus widely used in gold mining in antiquity.
In Shetland, tobacco was not smoked when first introduced, but instead was ground up into snuff, and inhaled up the nose. Snuff-querns consisted of an upper and lower stone, fixed together by a central iron pivot. The quern was held on the user's lap, the eye of the quern was filled with dried tobacco leaves, and then the upper-stone was turned using the handle. The friction caused by the turning ground the leaves into a fine powder that built up around the edge of the lower-stone. Many snuff-querns had a small hole or cut made near the edge of the upper-stone, into which a pointed end of a lamb's horn was placed in order to turn the stone; an alternative to using a handle.
Incidental uses
Quern stones’ ubiquity led to ancillary uses. For example, DeBoer, in his review of the traditional gambling games of North American tribes, reports that one of the games involved bouncing a group of split canes off a quern.
Quern stones may have been used as improvised weapons, as mentioned in the Bible: "But a certain woman threw an upper-millstone on Abimelech's head, and crushed his skull." (Judg. 9:53 NRSV)
Manufacture of quern-stones
The best type of stone from which to manufacture quern-stones are igneous rocks such as basalt. These have naturally rough surfaces, but grains do not detach easily, so the material being ground does not become gritty. However, such rocks are not always available, meaning that quern-stones have been manufactured from a wide variety of rocks, including sandstone, quartzite and limestone. Quernmore Crag near Lancaster in England is named after the quarrying of millstone grit used to make quern stones in these parts. The names of the mountains of Whernside and Great Whernside in the Yorkshire Dales have the same origin.
Rutter was able to show, for the southern Levant, that basalt quern-stones were preferred to those manufactured from other rock types. Basalt quern-stones were therefore transported over long-distances, leading him to argue that, despite their everyday, utilitarian function, they were also used as a status symbol.
Research in Scotland has indicated that to a degree regional styles existed.
Evolution of quern-stones
Knocking stones were used in the preparation of small quantities of cereal, however the earliest forms of quern were the saddle and trough querns. The earliest quern so far discovered dates to and was found at Abu Hureyra, Syria. A later development was the rotary quern, which takes several forms.
Saddle quern
The saddle quern is produced by rocking or rolling the muller using parallel motions (i.e., pushing and pulling the handstone), which forms a shape looking like a saddle. These are the most ancient and widely used type of quern-stone and were superseded around the 5th to the 4th century BC by the more efficient rotary quern. The handstones for saddle querns are generally either roughly cylindrical (not unlike a rolling pin) and used with both hands, or rough hemispheres and used with one hand. This provides a crushing motion, not a grinding action and is more suitable for crushing malted grain. It is not easy to produce flour from a saddle quern with unmalted grain. The muller is also referred to as a 'rubber' or 'mouler'. Some pictorial examples may be found in the metate article.
Rotary quern
As the name implies, the rotary quern used circular motions to grind the material, meaning both the upper and lower quern were generally circular. The handstone of a rotary quern is much heavier than that of saddle quern and provides the necessary weight for the grinding of unmalted grain into flour. In some cases the grinding surfaces of the stones fit into each other, the upper stone being slightly concave and the lower one convex.
Beehive quern
In this type, the upper stone is hemispherical, or bun-shaped, with a central conical hopper to hold the grain that falls down a hole to the grinding surface. It is held in position with a pivot that fits into a central hole in the bottom stone. The upper stone also has a deep horizontal socket in its steep side in which to place the wooden peg used as a handle to rotate or oscillate the upper stone. This was the earliest type of rotary quern to appear in the British Isles. It arrived in Britain in the middle of the Iron Age (about 400–300 BC) and spread into the northern half of Ireland, probably from Scotland, some time after the 2nd century BC.
Disc quern
Disc querns consist of two flat disc-shaped stones, an upper stone with a cylindrical perforation for containing a handle and a lower stone with a central spindle hole. The adjustable, discoid rotary quern has larger, flatter and more discoid stones than the beehive type. The lower stone was completely perforated. The long handle rotated in a shallow socket in the upper surface of the upper stone. They are thought to have originated in Spain 2,500 years ago and appear to have arrived in maritime Scotland from about 200 BC with people who built the defensive homes known as brochs. This Iron Age type closely resembles the adjustable Highland quern still in use in historic times. In Ireland, disc querns were the dominant type of querns in use between the years 500 CE–1500 CE, and usually had a diameter of 30–60 cm.
Garnett in his 1800 tour of Scotland describes the use of a hand quern as follows:
Miniature quern
Under in diameter and varying from roughly dressed to carefully worked, often with vertical handle sockets, a new class of querns has been identified having been overlooked in the past as weights, etc. In all respects they are like full sized quern stones and they show the typical wear signs that indicate that they were used for grinding small amounts of seeds, minerals or herbs. A suggestion that they may have been made as toys is thought to be unlikely.
Other types of quern stones
Other forms of quern-stone include hopper-rubbers and Pompeian mills, both used by the Romans. The larger rotary mills were usually worked by a donkey or horse via an extension arm of wood attached to the upper stone.
Querns utilizing crank-and-connecting rods were used in the Western Han Dynasty.
Laws against use
There was a legal requirement in Scotland for tenants to pay for use of the baron's mill. Early leases of mills gave to the miller the legal right to destroy quern-stones which were being used in defiance of thirlage agreements.
The obligations of thirlage eventually ceased to apply, but thirlage in Scotland was only formally and totally abolished on 28 November (Martinmas) 2004 by the Abolition of Feudal Tenure etc. (Scotland) Act 2000.
A similar obligation (the manorial right of mulcture) existed in England, established by custom only and not by statute. This was actively enforced using common law provisions. Among many recorded adjudications, the 1274 case of the abbot of St Albans Abbey, Cirencester, who owned as lord of the manor the exclusive milling rights for the whole town, is notable: he instructed his bailiffs to search townspeople's homes and seize any quern-stones they found. Eventually the town brought a case against the abbott at the town assizes. They lost, and were ordered to pay the abbot 100 marks in compensation for defaming him. To celebrate his victory, the abbot then used the stones to pave the floor of his private parlour.
Exercising the prerogative fell into gradual disuse and the last instance was in Wakefield, Yorkshire, when in 1850 the town bought out the owner of Soke Mill for .
Ornament and inscription
A number of handstones have been found with extra carving, however it is not always straightforward to separate decoration from practical functional purposes. The designs invest in the appearance of the handstone when it is in circular motion, and the ability of the quern-stones to change seeds into flour may invoke a feeling of transformative magic that attracted both reverence and status to these household objects. Three beehive querns found in Ireland have inscribed ornament of La Tène type, as do examples from England and Wales. Many of the horizontal slot-handled quern-stones have decoration that usually follows the basic pattern of motifs that encircle the hopper and/or the handle slot. One type though has an irregular pattern of cup marks that encircle the hopper.
A quern was discovered at Dunadd in Scotland that has a cross carved into the upper stone. The cross has expanded terminals and ultimately derives its form from Roman and Byzantine predecessors of the fifth and sixth centuries. This example has a high quality of finishing which reflects its 'cost' and enhances its symbolic value and social significance. The cross is likely to have 'protected' the corn and the resultant flour from evil, such as fungal rust or ergot. Various legends give miraculous power to mill-stones and several have been found which have been re-used in the construction of burial cists or as tomb stones. The association between quern stones and burial may be because they are used in the process of making bread, the staple of life. A broken or disused quern therefore can be seen as symbolic of death. In Clonmacnoise, near Athlone in County Offaly in Ireland, a quern stone was found which had been made into a tombstone, having been ornamented and the name Sechnasach, who died in 928 AD, inscribed onto it. A large quern was discovered on the Lough Scur crannog in Ireland.
The Wondrous Mauchline Quern
In the 9th century the Welsh monk Nennius wrote a history of Britain, the Historia Brittonum, in which he lists the thirteen wonders of Britain, and included in it is the wondrous 'Mauchline Quern' that ground constantly, except on Sundays. It could be heard working underground and the local placename 'Auchenbrain' may celebrate it, translating from the Gaelic as 'field of the quern'.
Related grinding tools
Knocking stone
Jato: a type of rotary quern-stone used in the Himalayan region of Nepal, Sikkim, Darjeeling and Bhutan.
(Malay)
Manos and metate
Millstone
| Technology | Industrial machinery | null |
924328 | https://en.wikipedia.org/wiki/Mole%20cricket | Mole cricket | Mole crickets are members of the insect family Gryllotalpidae, in the order Orthoptera (grasshoppers, locusts, and crickets). Mole crickets are cylindrical-bodied, fossorial insects about long as adults, with small eyes and shovel-like fore limbs highly developed for burrowing. They are present in many parts of the world and where they have arrived in new regions, may become agricultural pests.
Mole crickets have three life stages: eggs, nymphs, and adults. Most of their lives in these stages are spent underground, but adults have wings and disperse in the breeding season. They vary in their diet: some species are herbivores, mainly feeding on roots; others are omnivores, including worms and grubs in their diet; and a few are largely predatory. Male mole crickets have an exceptionally loud song; they sing from a burrow that opens out into the air in the shape of an exponential horn. The song is an almost pure tone, modulated into chirps. It is used to attract females, either for mating, or for indicating favourable habitats for them to lay their eggs.
In Zambia, mole crickets are thought to bring good fortune, while in Latin America, they are said to predict rain. In Florida, where Neoscapteriscus mole crickets are not native, they are considered pests, and various biological controls have been used. Gryllotalpa species have been used as food in West Java, Vietnam, Thailand, Laos, and the Philippines.
Description
Mole crickets vary in size and appearance, but most of them are of moderate size for an insect, typically between long as adults. They are adapted for underground life and are cylindrical in shape and covered with fine, dense hairs. The head, fore limbs, and prothorax are heavily sclerotised, but the abdomen is rather soft. The head bears two threadlike antennae and a pair of beady eyes. The two pairs of wings are folded flat over the abdomen; in most species, the fore wings are short and rounded and the hind wings are membranous and reach or exceed the tip of the abdomen; however, in some species, the hind wings are reduced in size and the insect is unable to fly. The fore legs are flattened for digging, but the hind legs are shaped somewhat like the legs of a true cricket; however, these limbs are more adapted for pushing soil, rather than leaping, which they do rarely and poorly. The nymphs resemble the adults apart from the absence of wings and genitalia; the wing pads become larger after each successive moult.
Taxonomy and phylogeny
The Gryllotalpidae are a monophyletic group in the order Orthoptera (grasshoppers, locusts, and crickets). Cladistic analysis of mole cricket morphology in 2015 identified six tribes, of which four were then new: Indioscaptorini (Scapteriscinae), Triamescaptorini, Gryllotalpellini and Neocurtillini (Gryllotalpinae), and two existing tribes, Scapteriscini and Gryllotalpini, are revised.
The group name is derived straightforwardly from Latin gryllus, cricket, and talpa, mole.
Within the extant subfamilies, genera include:
Gryllotalpinae
tribe Gryllotalpellini Cadena-Castañeda, 2015 (South America)
Gryllotalpella
tribe Gryllotalpini Leach, 1815 (World-wide)
Gryllotalpa
tribe Neocurtillini Cadena-Castañeda, 2015 (Americas)
Leptocurtilla
Neocurtilla
tribe Triamescaptorini Cadena-Castañeda, 2015 (New Zealand)
Triamescaptor
tribe incertae sedis
†genus Pterotriamescaptor
†genus Burmagryllotalpa Burmese amber, Myanmar, Cenomanian
Scapteriscinae
tribe Indioscaptorini Cadena-Castañeda, 2015 (Indian Subcontinent)
Indioscaptor
tribe Scapteriscini Zeuner, 1939 (Americas)
Scapteriscus
Neoscapteriscus
†subfamily Marchandiinae
†Archaeogryllotalpoides Martins-Neto 1991 Crato Formation, Brazil, Aptian
†Cratotetraspinus Martins-Neto 1997 Crato Formation, Brazil, Aptian
†Marchandia Perrichot et al. 2002 Charentese amber, France, Cenomanian
†Palaeoscapteriscops Martins-Neto 1991 Crato Formation, Brazil, Aptian
Incertae sedis
†Tresdigitus rectanguli Xu, Fang & Wang, 2020 Burmese amber, Myanmar, Cenomanian
Mole cricket fossils are rare. A stem group fossil, Cratotetraspinus, is known from the Lower Cretaceous of Brazil. Two specimens of Marchandia magnifica in amber have been found in the Lower Cretaceous of Charente-Maritime in France. They are somewhat more abundant in the Tertiary amber of the Baltic and Dominican regions; impressions are found in Europe and the American Green River Formation.
Convergent evolution
Mole crickets are not closely related to the "pygmy mole crickets", the Tridactyloidea, which are in the grasshopper suborder Caelifera rather than the cricket suborder Ensifera. The two groups, and indeed their resemblance in form to the mammalian mole family Talpidae with their powerful front limbs, form an example of convergent evolution, both developing adaptations for burrowing.
Behavior
Adults of most species of mole cricket can fly powerfully, if not with agility, but males do so infrequently. The females typically take wing soon after sunset, and are attracted to areas where males are calling, which they do for about an hour after sunset. This may be to mate, or they may be influenced by the suitability of the habitat for egg-laying, as demonstrated by the number of males present and calling in the vicinity.
Life cycle
Mole crickets undergo incomplete metamorphosis; when nymphs hatch from eggs, they increasingly resemble the adult form as they grow and pass through a series of up to 10 moults. After mating, a period of 1–2 weeks may occur before the female starts laying eggs. She burrows into the soil to a depth of , ( has been seen in the laboratory), and lays a clutch of 25 to 60 eggs. Neoscapteriscus females then retire, sealing the entrance passage, but in Gryllotalpa and Neocurtilla species, the female has been observed to remain in an adjoining chamber to tend the clutch. Further clutches may follow over several months, according to species. Eggs must be laid in moist ground, and many nymphs die because of insufficient moisture in the soil. The eggs hatch in a few weeks, and as they grow, the nymphs consume a great deal of plant material either underground or on the surface.
The adults of some species of mole crickets may move as far as during the breeding season. Mole crickets are active most of the year, but overwinter as nymphs or adults in cooler climates, resuming activity in the spring.
Burrowing
Mole crickets live almost entirely below ground, digging tunnels of different kinds for the major functions of life, including feeding, escape from predators, attracting a mate (by singing), mating, and raising of young.
Their main tunnels are used for feeding and for escape; they can dig themselves under ground very rapidly, and can move along existing tunnels at high speed both forwards and backwards. Their digging technique is to force the soil to either side with their powerful, shovel-like fore limbs, which are broad, flattened, toothed, and heavily sclerotised (the cuticle is hardened and darkened).
Males attract mates by constructing specially shaped tunnels in which they sing. Mating takes place in the male's burrow; the male may widen a tunnel to make room for the female to mount, though in some species, mating is tail-to-tail. Females lay their eggs either in their normal burrows or in specially dug brood chambers, which are sealed when complete in the case of the genus Neoscapteriscus or not sealed in the case of genera Gryllotalpa and Neocurtilla.
Song
Male mole crickets sing by stridulating, always under ground. In Gryllotalpa gryllotalpa, the song is based on an almost pure tone at 3.5 kHz, loud enough to make the ground vibrate 20 cm all round the burrow; in fact, the song is unique in each species. In G. gryllotalpa, the burrow is somewhat roughly sculpted; in G. vineae, the burrow is smooth and carefully shaped, with no irregularities larger than 1 mm. In both species, the burrow has two openings at the soil surface; at the other end is a constriction, then a resonating bulb, and then an escape tunnel. A burrow is used for at least a week. The male positions himself head down with his head in the bulb, and his tail is near the fork in the tunnel.
Mole crickets stridulate like other crickets by scraping the rear edge of the left fore wing, which forms a plectrum, against the lower surface of the right fore wing, which has a ratchet-like series of asymmetric teeth; the more acute edges face backwards, as do those of the plectrum. The plectrum can move forward with little resistance, but moving it backwards makes it catch each tooth, setting up a vibration in both wings. The sound-producing stroke is the raising (levation) of the wings. The resulting song resembles the result of modulating a pure tone with a 66-Hz wave to form regular chirps. In G. vineae, the wing levator muscle, which weighs 50 mg, can deliver 3.5 milliwatts of mechanical power; G. gryllotalpa can deliver about 1 milliwatt. G. vineae produces an exceptionally loud song from half an hour after sunset, continuing for an hour; it can be heard up to 600 m away. At a distance of 1 m from the burrow, the sound has a mean power over the stridulation cycle up to 88 decibels; the loudest recorded peak power was about 92 decibels; at the mouths of the burrow, the sound reaches around 115 decibels. G. gryllotalpa can deliver a peak sound pressure of 72 decibels and a mean of about 66 decibels. The throat of the horn appears to be tuned (offering low inductive reactance), making the burrow radiate sound efficiently; the efficiency increases when the burrow is wet and absorbs less sound. Mole crickets are the only insects that construct a sound-producing apparatus. Given the known sensitivity of a cricket's hearing (60 decibels), a night-flying G. vineae female should be able to detect the male's song at a range of 30 m; this compares to about 5 m for a typical Gryllus cricket that does not construct a burrow.
The loudness of the song is correlated with the size of the male and the quality of the habitat, both indicators of male attractiveness. The loudest males may attract 20 females in one evening, while a quieter male may attract none. This behaviour enables acoustic trapping; females can be trapped in large numbers by broadcasting a male's song very loudly.
Food
Mole crickets vary in their diets; some like the tawny mole cricket are herbivores, others are omnivores, feeding on larvae, worms, roots, and grasses, and others like the southern mole cricket are mainly predacious. They leave their burrows at night to forage for leaves and stems, which they drag underground before consumption, as well as consuming roots underground.
Predators, parasites, and pathogens
Besides birds, toads, and insectivorous mammals, the predators of mole crickets include subterranean assassin bugs, wolf spiders, and various beetles. The South American nematode Steinernema scapterisci kills Neocapteriscus mole crickets by introducing bacteria into their bodies, causing an overwhelming infection. Steinernema neocurtillae is native to Florida and attacks native Neocurtilla hexadactyla mole crickets. Parasitoid wasps of the genus Larra (Hymenoptera: Crabronidae) attack mole crickets, the female laying an egg on the external surface of the mole cricket, and the larva developing externally on the mole cricket host. Ormia depleta (Diptera: Tachinidae) is a specialized parasitoid of mole crickets in the genus Neoscapteriscus; the fly's larvae hatch from eggs inside her abdomen; she is attracted by the call of the male mole cricket and deposits a larva or more on any mole cricket individual (just as many females as males) with which she comes in contact. Specialist predators of mole cricket eggs in China and Japan include the bombardier beetle Stenaptinus jessoensis, whereas in South America, they include the bombardier beetle Pheropsophus aequinoctialis (Coleoptera: Carabidae); the adult beetle lays eggs near the burrows of mole crickets, and the beetle larvae find their way to the egg chamber and eat the eggs. Fungal diseases can devastate mole cricket populations during winters with sudden rises of temperature and thaws. The fungus Beauveria bassiana can overwhelm adult mole crickets and several other fungal, microsporidian, and viral pathogens have been identified.
Mole crickets evade predators by living below ground, and vigorously burrowing if disturbed at the surface. As a last-ditch defence, they eject a foul-smelling brown liquid from their anal glands when captured; they can also bite.
Distribution
Mole crickets are relatively common, but because they are nocturnal and spend nearly all their lives under ground in extensive tunnel systems, they are rarely seen. They inhabit agricultural fields and grassy areas. They are present in every continent except Antarctica; by 2014, 107 species had been described and more species are likely to be discovered, especially in Asia. Neoscapteriscus didactylus is a pest species, originating in South America; it has spread to the West Indies and New South Wales in Australia. Gryllotalpa africana is a major pest in South Africa; other Gryllotalpa species are widely distributed in Europe, Asia, and Australia. They are native to Britain (as to Western Europe), but the former population of G. gryllotalpa may now be extinct in mainland Britain, surviving in the Channel Islands.
Invasive mole crickets and their biological control
Invasive species are those that cause harm in their newly occupied area, where biological control may be attempted. The first-detected invasive mole cricket species was Neoscapteriscus didactylus, a South American species reported as a pest in St. Vincent, West Indies, as early as 1837; by 1900, it was a major agricultural pest in Puerto Rico. It had probably slowly expanded its range northwards, island by island, from South America. The only biological control program against N. didactylus was in Puerto Rico, and it succeeded in establishing the parasitoid wasp Larra bicolor from Amazonian Brazil. In 2001, N. didactylus in Puerto Rico seemed to be a pest only in irrigated crops and turf. Small-scale experimental applications of the nematode Steinernema scapterisci were made in irrigated turf, but survival of the nematode was poor. Very much later, this same species was reported as a pest in Queensland, Australia, presumably arriving by ship or plane. The next-detected invasive species was in the late 19th century in Hawaii, probably by ship. It was named as Gryllotalpa africana, but much later as perhaps G. orientalis. It was not identified as Gryllotalpa krishnani until 2020. It attacked sugarcane and was targeted with Larra polita from the Philippines in 1925, apparently successfully.
The next detection was in Georgia, USA, and at that time was assumed to be N. didactylus from the West Indies. It was, in fact, three South American Neoscapteriscus species, N. abbreviatus, N . vicinus, and N. borellii, probably having arrived in ship ballast. They caused major problems for decades as they spread in the Southeastern USA.Scapteriscus mole cricket populations had built up since the early decades of the 20th century and damaged pastures, lawns, playing fields, and vegetable crops. From the late 1940s, chlordane had been the insecticide of choice to control them, but when chlordane was banned by the U.S. EPA in the 1970s, ranchers were left with no economic and effective control method. Especially to aid Florida ranchers, a project that became known as the UF/IFAS Mole Cricket Research Program was initiated in 1978. In 1985, a multiple-authored report was published on accomplishments. In 1988, an account was published on prospects for biological control, and in 1996 an account of promising results with biological control. The program ended in 2004 after 25 years of running monitoring stations, and in 2006 a summary publication announced success: a 95% reduction in mole cricket numbers in northern Florida, with biological control agents spreading potentially to all parts of Florida. Efforts to use Larra bicolor as a biological control agent in Florida began by importing a stock from Puerto Rico. It became established in a small area of southeastern Florida, but had little effect on Neoscapteriscus populations. A stock from Bolivia became established in northern Florida and spread widely (with some help) to most of the rest of the state and neighboring states. Its survival depends upon the availability of suitable nectar sources.
Once gravid female Ormia depleta flies were found to be attracted to the song of Neoscapteriscus males in South America, a path to trap these flies with synthetic mole cricket songs was opened. Experimentation then led to a rearing method. Laborious rearing of over 10,000 flies on mole cricket hosts allowed releases of living fly pupae at many sites in Florida from the far northwest to the far south, mainly on golf courses, and mainly in 1989-1991. Populations were established, began to spread, and were monitored by use of synthetic mole cricket song. Eventually, the flies were found to have a continuous population from about 29°N then south to Miami, but the flies failed to survive the winter north of about 29°N. Shipment and release of the flies to states north of Florida was thus a wasted effort. As the flies had been imported from 23°S in Brazil and could not overwinter north of 29°N, whether flies from 30°S in Brazil might survive better in northern Florida was investigated in 1999, but they did not.
The third biological control agent to target Neoscapteriscus in Florida was the South American nematode Steinernema scapterisci. Small-scale releases proved it could persist for years in mole cricket-infested sandy Florida soils. Its use as a biopesticide against Neoscapteriscus was patented, making it attractive to industry. Industrial-scale production on artificial diet allowed large-scale trial applications in pastures and on golf courses, which succeeded in establishing populations in several counties, and these populations spread, but sales were disappointing, and the product was withdrawn from the market in 2014. Although experimental application was made in states north of Florida, only in southern Georgia was establishment of the nematode verified, suggesting little interest in the other states.
As pests
The main damage done by mole crickets is as a result of their burrowing activities. As they tunnel through the top few centimetres of soil, they push the ground up in little ridges, increasing evaporation of surface moisture, disturbing germinating seeds, and damaging the delicate young roots of seedlings. They are also injurious to turf and pasture grasses as they feed on their roots, leaving the plants prone to drying out and damage by use.
In their native lands, mole crickets have natural enemies that keep them under control. This is not the case when they have been accidentally introduced to other parts of the world. In Florida from the 1940s through the 1980s, they were considered pests and were described as "a serious problem". Their population densities have since declined greatly. A University of Florida entomology report suggests that South American Neoscapteriscus mole crickets may have entered the United States at Brunswick, Georgia, in ship's ballast from southern South America around 1899, but were at that time mistakenly believed to be from the West Indies. One possible remedy was biological pest control using the parasitoidal wasps Larra bicolor. Another remedy that has been successfully applied is use of the parasitic nematode Steinernema scapterisci. When this is applied in strips across grassland, it spreads throughout the pasture (and potentially beyond) within a few months and not only controls the mole crickets, but also remains infective in the soil for future years.
In human culture
Folklore and Literature
In Zambia, Gryllotalpa africana is held to bring good fortune to anyone who sees it.
In Uganda, they are believed to announce the death of a family member or a loved one. In ancient days, if it was found in the house, the elders advised the people to kill the bugs whilst crying so that it dies with the sorrow it had brought to the family.
In Latin America, Scapteriscus and Neocurtilla mole crickets are said to predict rain when they dig into the ground.
In Japan, in the past, they seem to have been associated with the worms/corpses/bugs that announce a person's sins to heaven in the Koshin or Koushin belief.
In ancient a Chinese poetry by Zhuang Zhou the poet’s own death is depicted in chapter 32 and features mole crickets:
As food
Gryllotalpa mole crickets have sometimes been used as food in West Java and Vietnam. In Thailand, mole crickets () are valued as food in Isan. They are usually eaten fried along with sticky rice.
In the Philippines, they are served as a delicacy called camaro in the province of Pampanga and are a tourist attraction. They are also served in parts of Northern Luzon. "Ararawan" or "susuhong" in Ilocano language is a delicacy or exotic food than june beetle or salagubang. It is found in Vintar, Ilocos Norte and is sold at P900 per kilo.
| Biology and health sciences | Orthoptera | null |
924422 | https://en.wikipedia.org/wiki/Hump-winged%20grig | Hump-winged grig | Hump-winged grigs are insects belonging to the genus Cyphoderris, in the family Prophalangopsidae, and superfamily Grylloidea (crickets). In modern times they are known only in northwestern North America and central Asia, but the fossil record indicates a wider distribution in the past.
There are three species in North America:
Cyphoderris buckelli - Buckell's grig
Cyphoderris monstrosa - Great grig
Cyphoderris strepitans - Sagebrush grig
Hump-winged grigs are known for their unique mating habits. Males call at night by sitting on a tree trunk with their head down and emitting a short, high-pitched trill. When a female mounts the male, the male uses two hooks on its back to hold onto the underside of the female's abdomen while transferring spermatophore. During copulation, the female eats the male's hind wings and drinks the male's blood for energy, causing permanent but nonfatal damage to the male. Hungry females are more likely to mate, will mount males sooner, and are less selective when choosing mating partners. "Virgin" males, with no hind wing damage, are generally more successful at mating than non-virgin males.
| Biology and health sciences | Orthoptera | Animals |
925570 | https://en.wikipedia.org/wiki/Skutterudite | Skutterudite | Skutterudite is a cobalt arsenide mineral containing variable amounts of nickel and iron substituting for cobalt with the ideal formula CoAs3. Some references give the arsenic a variable formula subscript of 2–3. High nickel varieties are referred to as nickel-skutterudite, previously chloanthite. It is a hydrothermal ore mineral found in moderate to high temperature veins with other Ni-Co minerals. Associated minerals are arsenopyrite, native silver, erythrite, annabergite, nickeline, cobaltite, silver sulfosalts, native bismuth, calcite, siderite, barite and quartz. It is mined as an ore of cobalt and nickel with a by-product of arsenic.
The crystal structure of this mineral has been found to be exhibited by several compounds with important technological uses.
The mineral has a bright metallic luster, and is tin white or light steel gray in color with a black streak. The specific gravity is 6.5 and the hardness is 5.5–6. Its crystal structure is isometric with cube and octahedron forms similar to that of pyrite. The arsenic content gives a garlic odor when heated or crushed.
Skutterudite was discovered in Skuterud Mines, Modum, Buskerud, Norway, in 1845. Smaltite is an alternative name for the mineral. Notable occurrences include Cobalt, Ontario, Skuterud, Norway, and Franklin, New Jersey, in the United States. The rare arsenide minerals are classified in Dana's sulfide mineral group, even though it contains no sulfur.
Crystal structure
The crystal structure of the skutterudite mineral was determined in 1928 by Oftedahl to be cubic, belonging to space group Im-3 (number 204). The unit cell can be considered to consist of eight smaller cubes made up of the Co atoms. Six of these cubes are filled with (almost) square planar rings of As, each of which is oriented parallel to one of the unit cell edges. The As atoms then form octahedra with Co in the center.
In crystallographic terms, the Co atoms occupy the 8c sites, while the As atoms occupy the 24g sites. The position of the Co atoms within the unit cell is fixed, while the positions of the As atoms are determined by the parameters x and y. It has been shown that for the As-rings to be fully square, these parameters must satisfy the Oftedahl relation x + y = 1/2. Any deviation from this relation yields a rectangular configuration of the As atoms; indeed, this is the case for all known compounds with this structure, and the As atoms do not form a perfect octahedron.
Together with the unit cell size and the assigned space group, the aforementioned parameters fully describe the crystal structure of the material. This structure is often referred to as the skutterudite structure.
Applications
Materials with a skutterudite structure are studied as a low cost thermoelectric material with low thermal conductivity. | Physical sciences | Minerals | Earth science |
925638 | https://en.wikipedia.org/wiki/Forsterite | Forsterite | Forsterite (Mg2SiO4; commonly abbreviated as Fo; also known as white olivine) is the magnesium-rich end-member of the olivine solid solution series. It is isomorphous with the iron-rich end-member, fayalite. Forsterite crystallizes in the orthorhombic system (space group Pbnm) with cell parameters a 4.75 Å (0.475 nm), b 10.20 Å (1.020 nm) and c 5.98 Å (0.598 nm).
Forsterite is associated with igneous and metamorphic rocks and has also been found in meteorites. In 2005 it was also found in cometary dust returned by the Stardust probe. In 2011 it was observed as tiny crystals in the dusty clouds of gas around a forming star.
Two polymorphs of forsterite are known: wadsleyite (also orthorhombic) and ringwoodite (isometric, cubic crystal system). Both are mainly known from meteorites.
Peridot is the gemstone variety of forsterite olivine.
Composition
Pure forsterite is composed of magnesium, oxygen and silicon. The chemical formula is Mg2SiO4. Forsterite, fayalite (Fe2SiO4) and tephroite (Mn2SiO4) are the end-members of the olivine solid solution series; other elements such as Ni and Ca substitute for Fe and Mg in olivine, but only in minor proportions in natural occurrences. Other minerals such as monticellite (CaMgSiO4), an uncommon calcium-rich mineral, share the olivine structure, but solid solution between olivine and these other minerals is limited. Monticellite is found in contact metamorphosed dolomites.
Geologic occurrence
Forsterite-rich olivine is the most abundant mineral in the mantle above a depth of about ; pyroxenes are also important minerals in this upper part of the mantle. Although pure forsterite does not occur in igneous rocks, dunite often contains olivine with forsterite contents at least as Mg-rich as (92% forsterite – 8% fayalite); common peridotite contains olivine typically at least as Mg-rich as . Due to its high melting point, olivine crystals are the first minerals to precipitate from a magmatic melt in a cumulate process, often with orthopyroxenes. Forsterite-rich olivine is a common crystallization product of mantle-derived magma. Olivine in mafic and ultramafic rocks typically is rich in the forsterite end-member.
Forsterite also occurs in dolomitic marble which results from the metamorphism of high magnesium limestones and dolomites.
Nearly pure forsterite occurs in some metamorphosed serpentinites. Fayalite-rich olivine is much less common. Nearly pure fayalite is a minor constituent in some granite-like rocks, and it is a major constituent of some metamorphic banded iron formations.
Structure, formation, and physical properties
Forsterite is mainly composed of the anion SiO44− and the cation Mg2+ in a molar ratio 1:2. Silicon is the central atom in the SiO44− anion. Each oxygen atom is bonded to the silicon by a single covalent bond. The four oxygen atoms have a partial negative charge because of the covalent bond with silicon. Therefore, oxygen atoms need to stay far from each other in order to reduce the repulsive force between them. The best geometry to reduce the repulsion is a tetrahedral shape. The cations occupy two different octahedral sites which are M1 and M2 and form ionic bonds with the silicate anions. M1 and M2 are slightly different. M2 site is larger and more regular than M1 as shown in Fig. 1. The packing in forsterite structure is dense. The space group of this structure is Pbnm and the point group is 2/m 2/m 2/m which is an orthorhombic crystal structure.
This structure of forsterite can form a complete solid solution by replacing the magnesium with iron. Iron can form two different cations which are Fe2+ and Fe3+. The iron(II) ion has the same charge as magnesium ion and it has a very similar ionic radius to magnesium. Consequently, Fe2+ can replace the magnesium ion in the olivine structure.
One of the important factors that can increase the portion of forsterite in the olivine solid solution is the ratio of iron(II) ions to iron(III) ions in the magma. As the iron(II) ions oxidize and become iron(III) ions, iron(III) ions cannot form olivine because of their 3+ charge. The occurrence of forsterite due to the oxidation of iron was observed in the Stromboli volcano in Italy. As the volcano fractured, gases and volatiles escaped from the magma chamber. The crystallization temperature of the magma increased as the gases escaped. Because iron(II) ions were oxidized in the Stromboli magma, little iron(II) was available to form Fe-rich olivine (fayalite). Hence, the crystallizing olivine was Mg-rich, and igneous rocks rich in forsterite were formed.
At high pressure, forsterite undergoes a phase transition into wadsleyite; under the conditions prevailing in the Earth's upper mantle, this transformation would occur at pressures of ca. 14–15 GPa. In high-pressure experiments, the transformation may be delayed so that forsterite can remain metastable at pressures up to almost 50 GPa (see fig.).
The progressive metamorphism between dolomite and quartz react to form forsterite, calcite and carbon dioxide:
2CaMg(CO3)2 + SiO2 -> Mg2SiO4 + 2CaCO3 + 2CO2
Forsterite reacts with quartz to form the orthopyroxene mineral enstatite in the following reaction:
Mg2SiO4 + SiO2 -> 2MgSiO3
Discovery and name
Forsterite was first described in 1824 for an occurrence at Mount Somma, Vesuvius, Italy. It was named by Armand Lévy in 1824 after the English naturalist and mineral collector Adolarius Jacob Forster.
Applications
Forsterite is being currently studied as a potential biomaterial for implants owing to its superior mechanical properties.
| Physical sciences | Silicate minerals | Earth science |
33051527 | https://en.wikipedia.org/wiki/Cell%20membrane | Cell membrane | The cell membrane (also known as the plasma membrane or cytoplasmic membrane, and historically referred to as the plasmalemma) is a biological membrane that separates and protects the interior of a cell from the outside environment (the extracellular space). The cell membrane consists of a lipid bilayer, made up of two layers of phospholipids with cholesterols (a lipid component) interspersed between them, maintaining appropriate membrane fluidity at various temperatures. The membrane also contains membrane proteins, including integral proteins that span the membrane and serve as membrane transporters, and peripheral proteins that loosely attach to the outer (peripheral) side of the cell membrane, acting as enzymes to facilitate interaction with the cell's environment. Glycolipids embedded in the outer lipid layer serve a similar purpose.
The cell membrane controls the movement of substances in and out of a cell, being selectively permeable to ions and organic molecules. In addition, cell membranes are involved in a variety of cellular processes such as cell adhesion, ion conductivity, and cell signalling and serve as the attachment surface for several extracellular structures, including the cell wall and the carbohydrate layer called the glycocalyx, as well as the intracellular network of protein fibers called the cytoskeleton. In the field of synthetic biology, cell membranes can be artificially reassembled.
History
Robert Hooke's discovery of cells in 1665 led to the proposal of the cell theory. Initially it was believed that all cells contained a hard cell wall since only plant cells could be observed at the time. Microscopists focused on the cell wall for well over 150 years until advances in microscopy were made. In the early 19th century, cells were recognized as being separate entities, unconnected, and bound by individual cell walls after it was found that plant cells could be separated. This theory extended to include animal cells to suggest a universal mechanism for cell protection and development.
By the second half of the 19th century, microscopy was still not advanced enough to make a distinction between cell membranes and cell walls. However, some microscopists correctly identified at this time that while invisible, it could be inferred that cell membranes existed in animal cells due to intracellular movement of components internally but not externally and that membranes were not the equivalent of a plant cell wall. It was also inferred that cell membranes were not vital components to all cells. Many refuted the existence of a cell membrane still towards the end of the 19th century. In 1890, a revision to the cell theory stated that cell membranes existed, but were merely secondary structures. It was not until later studies with osmosis and permeability that cell membranes gained more recognition. In 1895, Ernest Overton proposed that cell membranes were made of lipids.
The lipid bilayer hypothesis, proposed in 1925 by Gorter and Grendel, created speculation in the description of the cell membrane bilayer structure based on crystallographic studies and soap bubble observations. In an attempt to accept or reject the hypothesis, researchers measured membrane thickness. These researchers extracted the lipid from human red blood cells and measured the amount of surface area the lipid would cover when spread over the surface of the water. Since mature mammalian red blood cells lack both nuclei and cytoplasmic organelles, the plasma membrane is the only lipid-containing structure in the cell. Consequently, all of the lipids extracted from the cells can be assumed to have resided in the cells' plasma membranes. The ratio of the surface area of water covered by the extracted lipid to the surface area calculated for the red blood cells from which the lipid was 2:1(approx) and they concluded that the plasma membrane contains a lipid bilayer.
In 1925 it was determined by Fricke that the thickness of erythrocyte and yeast cell membranes ranged between 3.3 and 4 nm, a thickness compatible with a lipid monolayer. The choice of the dielectric constant used in these studies was called into question but future tests could not disprove the results of the initial experiment. Independently, the leptoscope was invented in order to measure very thin membranes by comparing the intensity of light reflected from a sample to the intensity of a membrane standard of known thickness. The instrument could resolve thicknesses that depended on pH measurements and the presence of membrane proteins that ranged from 8.6 to 23.2 nm, with the lower measurements supporting the lipid bilayer hypothesis. Later in the 1930s, the membrane structure model developed in general agreement to be the paucimolecular model of Davson and Danielli (1935). This model was based on studies of surface tension between oils and echinoderm eggs. Since the surface tension values appeared to be much lower than would be expected for an oil–water interface, it was assumed that some substance was responsible for lowering the interfacial tensions in the surface of cells. It was suggested that a lipid bilayer was in between two thin protein layers. The paucimolecular model immediately became popular and it dominated cell membrane studies for the following 30 years, until it became rivaled by the fluid mosaic model of Singer and Nicolson (1972).
Despite the numerous models of the cell membrane proposed prior to the fluid mosaic model, it remains the primary archetype for the cell membrane long after its inception in the 1970s. Although the fluid mosaic model has been modernized to detail contemporary discoveries, the basics have remained constant: the membrane is a lipid bilayer composed of hydrophilic exterior heads and a hydrophobic interior where proteins can interact with hydrophilic heads through polar interactions, but proteins that span the bilayer fully or partially have hydrophobic amino acids that interact with the non-polar lipid interior. The fluid mosaic model not only provided an accurate representation of membrane mechanics, it enhanced the study of hydrophobic forces, which would later develop into an essential descriptive limitation to describe biological macromolecules.
For many centuries, the scientists cited disagreed with the significance of the structure they were seeing as the cell membrane. For almost two centuries, the membranes were seen but mostly disregarded as an important structure with cellular function. It was not until the 20th century that the significance of the cell membrane as it was acknowledged. Finally, two scientists Gorter and Grendel (1925) made the discovery that the membrane is "lipid-based". From this, they furthered the idea that this structure would have to be in a formation that mimicked layers. Once studied further, it was found by comparing the sum of the cell surfaces and the surfaces of the lipids, a 2:1 ratio was estimated; thus, providing the first basis of the bilayer structure known today. This discovery initiated many new studies that arose globally within various fields of scientific studies, confirming that the structure and functions of the cell membrane are widely accepted.
The structure has been variously referred to by different writers as the ectoplast (de Vries, 1885), Plasmahaut (plasma skin, Pfeffer, 1877, 1891), Hautschicht (skin layer, Pfeffer, 1886; used with a different meaning by Hofmeister, 1867), plasmatic membrane (Pfeffer, 1900), plasma membrane, cytoplasmic membrane, cell envelope and cell membrane. Some authors who did not believe that there was a functional permeable boundary at the surface of the cell preferred to use the term plasmalemma (coined by Mast, 1924) for the external region of the cell.
Composition
Cell membranes contain a variety of biological molecules, notably lipids and proteins. Composition is not set, but constantly changing for fluidity and changes in the environment, even fluctuating during different stages of cell development. Specifically, the amount of cholesterol in human primary neuron cell membrane changes, and this change in composition affects fluidity throughout development stages.
Material is incorporated into the membrane, or deleted from it, by a variety of mechanisms:
Fusion of intracellular vesicles with the membrane (exocytosis) not only excretes the contents of the vesicle but also incorporates the vesicle membrane's components into the cell membrane. The membrane may form blebs around extracellular material that pinch off to become vesicles (endocytosis).
If a membrane is continuous with a tubular structure made of membrane material, then material from the tube can be drawn into the membrane continuously.
Although the concentration of membrane components in the aqueous phase is low (stable membrane components have low solubility in water), there is an exchange of molecules between the lipid and aqueous phases.
Lipids
The cell membrane consists of three classes of amphipathic lipids: phospholipids, glycolipids, and sterols. The amount of each depends upon the type of cell, but in the majority of cases phospholipids are the most abundant, often contributing for over 50% of all lipids in plasma membranes. Glycolipids only account for a minute amount of about 2% and sterols make up the rest. In red blood cell studies, 30% of the plasma membrane is lipid. However, for the majority of eukaryotic cells, the composition of plasma membranes is about half lipids and half proteins by weight.
The fatty chains in phospholipids and glycolipids usually contain an even number of carbon atoms, typically between 16 and 20. The 16- and 18-carbon fatty acids are the most common. Fatty acids may be saturated or unsaturated, with the configuration of the double bonds nearly always "cis". The length and the degree of unsaturation of fatty acid chains have a profound effect on membrane fluidity as unsaturated lipids create a kink, preventing the fatty acids from packing together as tightly, thus decreasing the melting temperature (increasing the fluidity) of the membrane. The ability of some organisms to regulate the fluidity of their cell membranes by altering lipid composition is called homeoviscous adaptation.
The entire membrane is held together via non-covalent interaction of hydrophobic tails, however the structure is quite fluid and not fixed rigidly in place. Under physiological conditions phospholipid molecules in the cell membrane are in the liquid crystalline state. It means the lipid molecules are free to diffuse and exhibit rapid lateral diffusion along the layer in which they are present. However, the exchange of phospholipid molecules between intracellular and extracellular leaflets of the bilayer is a very slow process. Lipid rafts and caveolae are examples of cholesterol-enriched microdomains in the cell membrane. Also, a fraction of the lipid in direct contact with integral membrane proteins, which is tightly bound to the protein surface is called annular lipid shell; it behaves as a part of protein complex.
Cholesterol is normally found dispersed in varying degrees throughout cell membranes, in the irregular spaces between the hydrophobic tails of the membrane lipids, where it confers a stiffening and strengthening effect on the membrane. Additionally, the amount of cholesterol in biological membranes varies between organisms, cell types, and even in individual cells. Cholesterol, a major component of plasma membranes, regulates the fluidity of the overall membrane, meaning that cholesterol controls the amount of movement of the various cell membrane components based on its concentrations. In high temperatures, cholesterol inhibits the movement of phospholipid fatty acid chains, causing a reduced permeability to small molecules and reduced membrane fluidity. The opposite is true for the role of cholesterol in cooler temperatures. Cholesterol production, and thus concentration, is up-regulated (increased) in response to cold temperature. At cold temperatures, cholesterol interferes with fatty acid chain interactions. Acting as antifreeze, cholesterol maintains the fluidity of the membrane. Cholesterol is more abundant in cold-weather animals than warm-weather animals. In plants, which lack cholesterol, related compounds called sterols perform the same function as cholesterol.
Phospholipids forming lipid vesicles
Lipid vesicles or liposomes are approximately spherical pockets that are enclosed by a lipid bilayer. These structures are used in laboratories to study the effects of chemicals in cells by delivering these chemicals directly to the cell, as well as getting more insight into cell membrane permeability. Lipid vesicles and liposomes are formed by first suspending a lipid in an aqueous solution then agitating the mixture through sonication, resulting in a vesicle. Measuring the rate of efflux from the inside of the vesicle to the ambient solution allows researchers to better understand membrane permeability. Vesicles can be formed with molecules and ions inside the vesicle by forming the vesicle with the desired molecule or ion present in the solution. Proteins can also be embedded into the membrane through solubilizing the desired proteins in the presence of detergents and attaching them to the phospholipids in which the liposome is formed. These provide researchers with a tool to examine various membrane protein functions.
Carbohydrates
Plasma membranes also contain carbohydrates, predominantly glycoproteins, but with some glycolipids (cerebrosides and gangliosides). Carbohydrates are important in the role of cell-cell recognition in eukaryotes; they are located on the surface of the cell where they recognize host cells and share information. Viruses that bind to cells using these receptors cause an infection. For the most part, no glycosylation occurs on membranes within the cell; rather generally glycosylation occurs on the extracellular surface of the plasma membrane. The glycocalyx is an important feature in all cells, especially epithelia with microvilli. Recent data suggest the glycocalyx participates in cell adhesion, lymphocyte homing, and many others. The penultimate sugar is galactose and the terminal sugar is sialic acid, as the sugar backbone is modified in the Golgi apparatus. Sialic acid carries a negative charge, providing an external barrier to charged particles.
Proteins
The cell membrane has large content of proteins, typically around 50% of membrane volume These proteins are important for the cell because they are responsible for various biological activities. Approximately a third of the genes in yeast code specifically for them, and this number is even higher in multicellular organisms. Membrane proteins consist of three main types: integral proteins, peripheral proteins, and lipid-anchored proteins.
As shown in the adjacent table, integral proteins are amphipathic transmembrane proteins. Examples of integral proteins include ion channels, proton pumps, and g-protein coupled receptors. Ion channels allow inorganic ions such as sodium, potassium, calcium, or chlorine to diffuse down their electrochemical gradient across the lipid bilayer through hydrophilic pores across the membrane. The electrical behavior of cells (i.e. nerve cells) is controlled by ion channels. Proton pumps are protein pumps that are embedded in the lipid bilayer that allow protons to travel through the membrane by transferring from one amino acid side chain to another. Processes such as electron transport and generating ATP use proton pumps. A G-protein coupled receptor is a single polypeptide chain that crosses the lipid bilayer seven times responding to signal molecules (i.e. hormones and neurotransmitters). G-protein coupled receptors are used in processes such as cell to cell signaling, the regulation of the production of cAMP, and the regulation of ion channels.
The cell membrane, being exposed to the outside environment, is an important site of cell–cell communication. As such, a large variety of protein receptors and identification proteins, such as antigens, are present on the surface of the membrane. Functions of membrane proteins can also include cell–cell contact, surface recognition, cytoskeleton contact, signaling, enzymatic activity, or transporting substances across the membrane.
Most membrane proteins must be inserted in some way into the membrane. For this to occur, an N-terminus "signal sequence" of amino acids directs proteins to the endoplasmic reticulum, which inserts the proteins into a lipid bilayer. Once inserted, the proteins are then transported to their final destination in vesicles, where the vesicle fuses with the target membrane.
Function
The cell membrane surrounds the cytoplasm of living cells, physically separating the intracellular components from the extracellular environment. The cell membrane also plays a role in anchoring the cytoskeleton to provide shape to the cell, and in attaching to the extracellular matrix and other cells to hold them together to form tissues. Fungi, bacteria, most archaea, and plants also have a cell wall, which provides a mechanical support to the cell and precludes the passage of larger molecules.
The cell membrane is selectively permeable and able to regulate what enters and exits the cell, thus facilitating the transport of materials needed for survival. The movement of substances across the membrane can be achieved by either passive transport, occurring without the input of cellular energy, or by active transport, requiring the cell to expend energy in transporting it. The membrane also maintains the cell potential. The cell membrane thus works as a selective filter that allows only certain things to come inside or go outside the cell. The cell employs a number of transport mechanisms that involve biological membranes:
Passive osmosis and diffusion: Some substances (small molecules, ions) such as carbon dioxide (CO2) and oxygen (O2), can move across the plasma membrane by diffusion, which is a passive transport process. Because the membrane acts as a barrier for certain molecules and ions, they can occur in different concentrations on the two sides of the membrane. Diffusion occurs when small molecules and ions move freely from high concentration to low concentration in order to equilibrate the membrane. It is considered a passive transport process because it does not require energy and is propelled by the concentration gradient created by each side of the membrane. Such a concentration gradient across a semipermeable membrane sets up an osmotic flow for the water. Osmosis, in biological systems involves a solvent, moving through a semipermeable membrane similarly to passive diffusion as the solvent still moves with the concentration gradient and requires no energy. While water is the most common solvent in cell, it can also be other liquids as well as supercritical liquids and gases.
Transmembrane protein channels and transporters: Transmembrane proteins extend through the lipid bilayer of the membranes; they function on both sides of the membrane to transport molecules across it. Nutrients, such as sugars or amino acids, must enter the cell, and certain products of metabolism must leave the cell. Such molecules can diffuse passively through protein channels such as aquaporins in facilitated diffusion or are pumped across the membrane by transmembrane transporters. Protein channel proteins, also called permeases, are usually quite specific, and they only recognize and transport a limited variety of chemical substances, often limited to a single substance. Another example of a transmembrane protein is a cell-surface receptor, which allow cell signaling molecules to communicate between cells.
Endocytosis: Endocytosis is the process in which cells absorb molecules by engulfing them. The plasma membrane creates a small deformation inward, called an invagination, in which the substance to be transported is captured. This invagination is caused by proteins on the outside on the cell membrane, acting as receptors and clustering into depressions that eventually promote accumulation of more proteins and lipids on the cytosolic side of the membrane. The deformation then pinches off from the membrane on the inside of the cell, creating a vesicle containing the captured substance. Endocytosis is a pathway for internalizing solid particles ("cell eating" or phagocytosis), small molecules and ions ("cell drinking" or pinocytosis), and macromolecules. Endocytosis requires energy and is thus a form of active transport.
Exocytosis: Just as material can be brought into the cell by invagination and formation of a vesicle, the membrane of a vesicle can be fused with the plasma membrane, extruding its contents to the surrounding medium. This is the process of exocytosis. Exocytosis occurs in various cells to remove undigested residues of substances brought in by endocytosis, to secrete substances such as hormones and enzymes, and to transport a substance completely across a cellular barrier. In the process of exocytosis, the undigested waste-containing food vacuole or the secretory vesicle budded from Golgi apparatus, is first moved by cytoskeleton from the interior of the cell to the surface. The vesicle membrane comes in contact with the plasma membrane. The lipid molecules of the two bilayers rearrange themselves and the two membranes are, thus, fused. A passage is formed in the fused membrane and the vesicles discharges its contents outside the cell.
Prokaryotes
Prokaryotes are divided into two different groups, Archaea and Bacteria, with bacteria dividing further into gram-positive and gram-negative. Gram-negative bacteria have both a plasma membrane and an outer membrane separated by periplasm; however, other prokaryotes have only a plasma membrane. These two membranes differ in many aspects. The outer membrane of the gram-negative bacteria differs from other prokaryotes due to phospholipids forming the exterior of the bilayer, and lipoproteins and phospholipids forming the interior. The outer membrane typically has a porous quality due to its presence of membrane proteins, such as gram-negative porins, which are pore-forming proteins. The inner plasma membrane is also generally symmetric whereas the outer membrane is asymmetric because of proteins such as the aforementioned.
Also, for the prokaryotic membranes, there are multiple things that can affect the fluidity. One of the major factors that can affect the fluidity is fatty acid composition. For example, when the bacteria Staphylococcus aureus was grown at 37 °C for 24 h, the membrane exhibited a more fluid state instead of a gel-like state. This supports the concept that in higher temperatures, the membrane is more fluid than in colder temperatures. When the membrane is becoming more fluid and needs to become more stabilized, it will make longer fatty acid chains or saturated fatty acid chains in order to help stabilize the membrane.
Bacteria are also surrounded by a cell wall composed of peptidoglycan (amino acids and sugars). Some eukaryotic cells also have cell walls, but none that are made of peptidoglycan. The outer membrane of gram negative bacteria is rich in lipopolysaccharides, which are combined poly- or oligosaccharide and carbohydrate lipid regions that stimulate the cell's natural immunity. The outer membrane can bleb out into periplasmic protrusions under stress conditions or upon virulence requirements while encountering a host target cell, and thus such blebs may work as virulence organelles. Bacterial cells provide numerous examples of the diverse ways in which prokaryotic cell membranes are adapted with structures that suit the organism's niche. For example, proteins on the surface of certain bacterial cells aid in their gliding motion. Many gram-negative bacteria have cell membranes which contain ATP-driven protein exporting systems.
Structures
Fluid mosaic model
According to the fluid mosaic model of S. J. Singer and G. L. Nicolson (1972), which replaced the earlier model of Davson and Danielli, biological membranes can be considered as a two-dimensional liquid in which lipid and protein molecules diffuse more or less easily. Although the lipid bilayers that form the basis of the membranes do indeed form two-dimensional liquids by themselves, the plasma membrane also contains a large quantity of proteins, which provide more structure. Examples of such structures are protein-protein complexes, pickets and fences formed by the actin-based cytoskeleton, and potentially lipid rafts.
Lipid bilayer
Lipid bilayers form through the process of self-assembly. The cell membrane consists primarily of a thin layer of amphipathic phospholipids that spontaneously arrange so that the hydrophobic "tail" regions are isolated from the surrounding water while the hydrophilic "head" regions interact with the intracellular (cytosolic) and extracellular faces of the resulting bilayer. This forms a continuous, spherical lipid bilayer. Hydrophobic interactions (also known as the hydrophobic effect) are the major driving forces in the formation of lipid bilayers. An increase in interactions between hydrophobic molecules (causing clustering of hydrophobic regions) allows water molecules to bond more freely with each other, increasing the entropy of the system. This complex interaction can include noncovalent interactions such as van der Waals, electrostatic and hydrogen bonds.
Lipid bilayers are generally impermeable to ions and polar molecules. The arrangement of hydrophilic heads and hydrophobic tails of the lipid bilayer prevent polar solutes (ex. amino acids, nucleic acids, carbohydrates, proteins, and ions) from diffusing across the membrane, but generally allows for the passive diffusion of hydrophobic molecules. This affords the cell the ability to control the movement of these substances via transmembrane protein complexes such as pores, channels and gates.
Flippases and scramblases concentrate phosphatidyl serine, which carries a negative charge, on the inner membrane. Along with NANA, this creates an extra barrier to charged moieties moving through the membrane.
Membranes serve diverse functions in eukaryotic and prokaryotic cells. One important role is to regulate the movement of materials into and out of cells. The phospholipid bilayer structure (fluid mosaic model) with specific membrane proteins accounts for the selective permeability of the membrane and passive and active transport mechanisms. In addition, membranes in prokaryotes and in the mitochondria and chloroplasts of eukaryotes facilitate the synthesis of ATP through chemiosmosis.
Membrane polarity
The apical membrane or luminal membrane of a polarized cell is the surface of the plasma membrane that faces inward to the lumen. This is particularly evident in epithelial and endothelial cells, but also describes other polarized cells, such as neurons. The basolateral membrane or basolateral cell membrane of a polarized cell is the surface of the plasma membrane that forms its basal and lateral surfaces. It faces outwards, towards the interstitium, and away from the lumen. Basolateral membrane is a compound phrase referring to the terms "basal (base) membrane" and "lateral (side) membrane", which, especially in epithelial cells, are identical in composition and activity. Proteins (such as ion channels and pumps) are free to move from the basal to the lateral surface of the cell or vice versa in accordance with the fluid mosaic model. Tight junctions join epithelial cells near their apical surface to prevent the migration of proteins from the basolateral membrane to the apical membrane. The basal and lateral surfaces thus remain roughly equivalent to one another, yet distinct from the apical surface.
Membrane structures
Cell membrane can form different types of "supramembrane" structures such as caveolae, postsynaptic density, podosomes, invadopodia, focal adhesion, and different types of cell junctions. These structures are usually responsible for cell adhesion, communication, endocytosis and exocytosis. They can be visualized by electron microscopy or fluorescence microscopy. They are composed of specific proteins, such as integrins and cadherins.
Cytoskeleton
The cytoskeleton is found underlying the cell membrane in the cytoplasm and provides a scaffolding for membrane proteins to anchor to, as well as forming organelles that extend from the cell. Indeed, cytoskeletal elements interact extensively and intimately with the cell membrane. Anchoring proteins restricts them to a particular cell surface — for example, the apical surface of epithelial cells that line the vertebrate gut — and limits how far they may diffuse within the bilayer. The cytoskeleton is able to form appendage-like organelles, such as cilia, which are microtubule-based extensions covered by the cell membrane, and filopodia, which are actin-based extensions. These extensions are ensheathed in membrane and project from the surface of the cell in order to sense the external environment and/or make contact with the substrate or other cells. The apical surfaces of epithelial cells are dense with actin-based finger-like projections known as microvilli, which increase cell surface area and thereby increase the absorption rate of nutrients. Localized decoupling of the cytoskeleton and cell membrane results in formation of a bleb.
Intracellular membranes
The content of the cell, inside the cell membrane, is composed of numerous membrane-bound organelles, which contribute to the overall function of the cell. The origin, structure, and function of each organelle leads to a large variation in the cell composition due to the individual uniqueness associated with each organelle.
Mitochondria and chloroplasts are considered to have evolved from bacteria, known as the endosymbiotic theory. This theory arose from the idea that Paracoccus and Rhodopseudomonas, types of bacteria, share similar functions to mitochondria and blue-green algae (cyanobacteria) share similar functions to chloroplasts. Endosymbiotic theory proposes that through the course of evolution, a eukaryotic cell engulfed these two types of bacteria, leading to the formation of mitochondria and chloroplasts inside eukaryotic cells. This engulfment lead to the double-membranes systems of these organelles in which the outer membrane originated from the host's plasma membrane and the inner membrane was the endosymbiont's plasma membrane. Considering that mitochondria and chloroplasts both contain their own DNA is further support that both of these organelles evolved from engulfed bacteria that thrived inside a eukaryotic cell.
In eukaryotic cells, the nuclear membrane separates the contents of the nucleus from the cytoplasm of the cell. The nuclear membrane is formed by an inner and outer membrane, providing the strict regulation of materials in to and out of the nucleus. Materials move between the cytosol and the nucleus through nuclear pores in the nuclear membrane. If a cell's nucleus is more active in transcription, its membrane will have more pores. The protein composition of the nucleus can vary greatly from the cytosol as many proteins are unable to cross through pores via diffusion. Within the nuclear membrane, the inner and outer membranes vary in protein composition, and only the outer membrane is continuous with the endoplasmic reticulum (ER) membrane. Like the ER, the outer membrane also possesses ribosomes responsible for producing and transporting proteins into the space between the two membranes. The nuclear membrane disassembles during the early stages of mitosis and reassembles in later stages of mitosis.
The ER, which is part of the endomembrane system, which makes up a very large portion of the cell's total membrane content. The ER is an enclosed network of tubules and sacs, and its main functions include protein synthesis, and lipid metabolism. There are 2 types of ER, smooth and rough. The rough ER has ribosomes attached to it used for protein synthesis, while the smooth ER is used more for the processing of toxins and calcium regulation in the cell.
The Golgi apparatus has two interconnected round Golgi cisternae. Compartments of the apparatus forms multiple tubular-reticular networks responsible for organization, stack connection and cargo transport that display a continuous grape-like stringed vesicles ranging from 50 to 60 nm. The apparatus consists of three main compartments, a flat disc-shaped cisterna with tubular-reticular networks and vesicles.
Variations
The cell membrane has different lipid and protein compositions in distinct types of cells and may have therefore specific names for certain cell types.
Sarcolemma in muscle cells: Sarcolemma is the name given to the cell membrane of muscle cells. Although the sarcolemma is similar to other cell membranes, it has other functions that set it apart. For instance, the sarcolemma transmits synaptic signals, helps generate action potentials, and is very involved in muscle contraction. Unlike other cell membranes, the sarcolemma makes up small channels called T-tubules that pass through the entirety of muscle cells. It has also been found that the average sarcolemma is 10 nm thick as opposed to the 4 nm thickness of a general cell membrane.
Oolemma is the cell membrane in oocytes: The oolemma of oocytes, (immature egg cells) are not consistent with a lipid bilayer as they lack a bilayer and do not consist of lipids. Rather, the structure has an inner layer, the fertilization envelope, and the exterior is made up of the vitelline layer, which is made up of glycoproteins; however, channels and proteins are still present for their functions in the membrane.
Axolemma: The specialized plasma membrane on the axons of nerve cells that is responsible for the generation of the action potential. It consists of a granular, densely packed lipid bilayer that works closely with the cytoskeleton components spectrin and actin. These cytoskeleton components are able to bind to and interact with transmembrane proteins in the axolemma.
Permeability
The permeability of a membrane is the rate of passive diffusion of molecules through the membrane. These molecules are known as permeant molecules. Permeability depends mainly on the electric charge and polarity of the molecule and to a lesser extent the molar mass of the molecule. Due to the cell membrane's hydrophobic nature, small electrically neutral molecules pass through the membrane more easily than charged, large ones. The inability of charged molecules to pass through the cell membrane results in pH partition of substances throughout the fluid compartments of the body.
| Biology and health sciences | Organelles and other cell parts | null |
3842182 | https://en.wikipedia.org/wiki/Cord%20%28unit%29 | Cord (unit) | The cord is a unit of measure of dry volume used to measure firewood and pulpwood in the United States and Canada.
A cord is the amount of wood that, when "racked and well stowed" (arranged so pieces are aligned, parallel, touching, and compact), occupies a volume of . This corresponds to a well-stacked woodpile high, wide, and deep; or any other arrangement of linear measurements that yields the same volume.
The name cord probably comes from the use of a cord or string to measure it.
The face cord is a unit of volume for stacked firewood, 8 feet long, 4 feet wide, and 16 inches high—equal to 1/3 of a cord. The symbol for the unit is fc - cd.
Definitions
In Canada, the cord is legally defined by Measurement Canada. The cord is one of three legal standards for the sale of firewood in Canada: stacked cubic meter (or stere), cubic foot, and cord.
In the United States, the cord is defined by statute in most states. The U.S. National Institute of Standards and Technology Handbook 130, section 2.4.1.2, defines a cord and provides uniform regulations for the sale of fireplace and stove wood. In the metric system, wood is usually measured in steres and cubic meters: 1 stere = 1 m3 ≈ 0.276 cords.
Maine appears unique among U.S. states by also defining a "loose thrown cord" or pile of cut firewood: "A cord of in length shall mean the amount of wood, bark, and air contained in a space of ; and a cord of wood in length shall mean the amount of wood, bark, and air contained in a space of . [1981, c. 219 (amd).]"
Other non-official terms for firewood volume include standing cord, kitchen cord, running cord, face cord, fencing cord, country cord, long cord, and rick, all subject to local variation. These are usually taken to mean a well-stacked pile of wood in which the logs are shorter or longer than in a legal cord, to accommodate various burners. For example, a face cord commonly consists of wood that is long. The volume of a face cord therefore is typically 1/3 of the volume of a full cord even though it is long and high. A face cord is also called a rick in the midwestern United States.
The term is used in other English-speaking countries, such as New Zealand, but may not have a legal definition.
The corde was a unit of volume used before metrication in several French-speaking countries (France, Belgium and Luxembourg). Its value varied from depending on the region, corresponding approximately to 2 to 5 steres.
Heating value
One seasoned (dry) cord of Northern red oak with a heating value of has the heating equivalent of of fuel oil with a heating value of .
Australia
Until metrication in Australia, an imperial cord was a measurement for wood and firewood. The measurements for a cord of wood were 4 feet wide by 4 feet deep by 8 feet high, or usually a stack of wood containing 128 cubic feet (cu ft). The imperial cord enclosed 128 cu ft.
| Physical sciences | Volume | Basics and measurement |
3842527 | https://en.wikipedia.org/wiki/Dry%20measure | Dry measure | Dry measures are units of volume to measure bulk commodities that are not fluids and that were typically shipped and sold in standardized containers such as barrels. They have largely been replaced by the units used for measuring volumes in the metric system and liquid volumes in the imperial system but are still used for some commodities in the US customary system. They were or are typically used in agriculture, agronomy, and commodity markets to measure grain, dried beans, dried and fresh produce, and some seafood. They were formerly used for many other foods, such as salt pork and salted fish, and for industrial commodities such as coal, cement, and lime.
The names are often the same as for the units used to measure liquids, despite representing different volumes. The larger volumes of the dry measures apparently arose because they were based on heaped rather than "struck" (leveled) containers.
Today, many units nominally of dry measure have become standardized as units of mass (see bushel); and many other units are commonly conflated or confused with units of mass.
Metric units
In the original metric system, the unit of dry volume was the stere, equal to a one-meter cube, but this is not part of the modern metric system; the liter and the cubic meter are now used. However, the stere is still widely used for firewood.
Imperial and US customary units
In US customary units, most units of volume exist both in a dry and a liquid version, with the same name, but different values: the dry hogshead, dry barrel, dry gallon, dry quart, dry pint, etc. The bushel and the peck are only used for dry goods. Imperial units of volume are the same for both dry and liquid goods. They have a different value from both the dry and liquid US versions.
Many of the units are associated with particular goods, so for instance the dry hogshead has been used for sugar and for tobacco, and the peck for apples. There are also special measures for specific goods, such as the cord of wood, the sack, the bale of wool or cotton, the box of fruit, etc.
Because it is difficult to measure actual volume and easy to measure mass, many of these units are now also defined as units of mass, specific to each commodity, so a bushel of apples is a different weight from a bushel of wheat (weighed at a specific moisture level). Indeed, the bushel, the best-known unit of dry measure because it is the quoted unit in commodity markets, is in fact a unit of mass in those contexts.
Conversely, the ton used in specifying tonnage and in freight calculations is often a volume measurement rather than a mass measurement.
In US cooking, dry and liquid measures are the same: the cup, the tablespoon, the teaspoon.
US dry measures are 16% larger than liquid measures.
Struck and heaped measurement
The volume of bulk goods is usually measured by filling a standard container, so the containers' names and the units' names are often the same, and indeed both are called "measures". Normally, a level or struck measure is assumed, with the excess being swept off level ("struck") with the measure's brim—the stick used for this is called a "strickle". Sometimes heaped or heaping measures are used, with the commodity heaped in a cone above the measure.
There was historically a tendency for landowners to demand heaped bushels of commodities from their peasants, while at the same time peasants were obliged to purchase commodities from stricken containers. Rules outlawing this practice were circumvented through use of heavy round strickles, which would compress the contents of a bushel.
US units of dry measure
(rounded to 4 digits)
| Physical sciences | Basics | Basics and measurement |
3842768 | https://en.wikipedia.org/wiki/Network%20on%20a%20chip | Network on a chip | A network on a chip or network-on-chip (NoC or ) is a network-based communications subsystem on an integrated circuit ("microchip"), most typically between modules in a system on a chip (SoC). The modules on the IC are typically semiconductor IP cores schematizing various functions of the computer system, and are designed to be modular in the sense of network science. The network on chip is a router-based packet switching network between SoC modules.
NoC technology applies the theory and methods of computer networking to on-chip communication and brings notable improvements over conventional bus and crossbar communication architectures. Networks-on-chip come in many network topologies, many of which are still experimental as of 2018.
In 2000s, researchers had started to propose a type of on-chip interconnection in the form of packet switching networks in order to address the scalability issues of bus-based design. Preceding researches proposed the design that routes data packets instead of routing the wires. Then, the concept of "network on chips" was proposed in 2002. NoCs improve the scalability of systems-on-chip and the power efficiency of complex SoCs compared to other communication subsystem designs. They are an emerging technology, with projections for large growth in the near future as multicore computer architectures become more common.
Structure
NoCs can span synchronous and asynchronous clock domains, known as clock domain crossing, or use unclocked asynchronous logic. NoCs support globally asynchronous, locally synchronous electronics architectures, allowing each processor core or functional unit on the System-on-Chip to have its own clock domain.
Architectures
NoC architectures typically model sparse small-world networks (SWNs) and scale-free networks (SFNs) to limit the number, length, area and power consumption of interconnection wires and point-to-point connections.
Topology
The topology determines the physical layout and connections between nodes and channels. The message traverses hops, and each hop's channel length depends on the topology. The topology significantly influences both latency and power consumption. Furthermore, since the topology determines the number of alternative paths between nodes, it affects the network traffic distribution, and hence the network bandwidth and performance achieved.
Benefits
Traditionally, ICs have been designed with dedicated point-to-point connections, with one wire dedicated to each signal. This results in a dense network topology. For large designs, in particular, this has several limitations from a physical design viewpoint. It requires power quadratic in the number of interconnections. The wires occupy much of the area of the chip, and in nanometer CMOS technology, interconnects dominate both performance and dynamic power dissipation, as signal propagation in wires across the chip requires multiple clock cycles. This also allows more parasitic capacitance, resistance and inductance to accrue on the circuit. (See Rent's rule for a discussion of wiring requirements for point-to-point connections).
Sparsity and locality of interconnections in the communications subsystem yield several improvements over traditional bus-based and crossbar-based systems.
Parallelism and scalability
The wires in the links of the network-on-chip are shared by many signals. A high level of parallelism is achieved, because all data links in the NoC can operate simultaneously on different data packets. Therefore, as the complexity of integrated systems keeps growing, a NoC provides enhanced performance (such as throughput) and scalability in comparison with previous communication architectures (e.g., dedicated point-to-point signal wires, shared buses, or segmented buses with bridges). The algorithms must be designed in such a way that they offer large parallelism and can hence utilize the potential of NoC.
Current research
Some researchers think that NoCs need to support quality of service (QoS), namely achieve the various requirements in terms of throughput, end-to-end delays, fairness, and deadlines. Real-time computation, including audio and video playback, is one reason for providing QoS support. However, current system implementations like VxWorks, RTLinux or QNX are able to achieve sub-millisecond real-time computing without special hardware.
This may indicate that for many real-time applications the service quality of existing on-chip interconnect infrastructure is sufficient, and dedicated hardware logic would be necessary to achieve microsecond precision, a degree that is rarely needed in practice for end users (sound or video jitter need only tenth of milliseconds latency guarantee). Another motivation for NoC-level quality of service (QoS) is to support multiple concurrent users sharing resources of a single chip multiprocessor in a public cloud computing infrastructure. In such instances, hardware QoS logic enables the service provider to make contractual guarantees on the level of service that a user receives, a feature that may be deemed desirable by some corporate or government clients.
Many challenging research problems remain to be solved at all levels, from the physical link level through the network level, and all the way up to the system architecture and application software. The first dedicated research symposium on networks on chip was held at Princeton University, in May 2007. The second IEEE International Symposium on Networks-on-Chip was held in April 2008 at Newcastle University.
Research has been conducted on integrated optical waveguides and devices comprising an optical network on a chip (ONoC).
The possible way to increasing the performance of NoC is use wireless communication channels between chiplets — named wireless network on chip (WiNoC).
Side benefits
In a multi-core system, connected by NoC, coherency messages and cache miss requests have to pass switches. Accordingly, switches can be augmented with simple tracking and forwarding elements to detect which cache blocks will be requested in the future by which cores. Then, the forwarding elements multicast any requested block to all the cores that may request the block in the future. This mechanism reduces cache miss rate.
Benchmarks
NoC development and studies require comparing different proposals and options. NoC traffic patterns are under development to help such evaluations. Existing NoC benchmarks include NoCBench and MCSL NoC Traffic Patterns.
Interconnect processing unit
An interconnect processing unit (IPU) is an on-chip communication network with hardware and software components which jointly implement key functions of different system-on-chip programming models through a set of communication and synchronization primitives and provide low-level platform services to enable advanced features in modern heterogeneous applications on a single die.
| Technology | Semiconductors | null |
24425329 | https://en.wikipedia.org/wiki/Belt%20Supergroup | Belt Supergroup | The Belt Supergroup is an assemblage of primarily fine-grained sedimentary rocks and mafic intrusive rocks of late Precambrian (Mesoproterozoic) age. It is more than thick, covers an area of some 200,000 km2 (77,220 sq. mi), and is considered to be one of the world's best-exposed and most accessible sequences of Mesoproterozoic rocks. It was named after the Big Belt Mountains in west-central Montana. It is present in western Montana and northern Idaho, with minor occurrences in northeastern Washington and western Wyoming. It extends into Canada where the equivalent rocks, which are called the Purcell Supergroup, are exposed in southeastern British Columbia and southwestern Alberta. The rocks of the Belt Supergroup contain economically significant deposits of lead, zinc, silver, copper, gold, and other metals in a number of areas, and some of the Belt rocks contain fossil stromatolites.
Spectacular outcrops of Belt rocks can be seen in Glacier National Park in northwestern Montana and in Waterton Lakes National Park in southwestern Alberta.
Lithology and sedimentology
The Belt Supergroup is dominated by fine-grained sedimentary rocks, primarily mudstones, siltstones, fine-grained quartzose sandstones and limestones. Most have undergone weak metamorphism to greenschist facies, and as a result the mudrocks are commonly classified as argillites and the sandstones as quartzites. The Belt Supergroup also includes lesser amounts of coarser grained sandstones and conglomerates. Mafic intrusive rocks are present locally in the lower portion.
Much of the sedimentation probably occurred between about 1450 and 1400 Ma (million years) ago. Sedimentary structures are well preserved in most of the Belt rocks despite their great age. The sedimentation is unusual in that 1) there is an abundance of fine-grained sediment and very little coarser sediment, 2) there is a lack of sequence boundaries that are common in Phanerozoic sediments, and 3) cyclic and rhythmic deposition occurred over long periods of time. The Belt Supergroup is also noted for "Molar Tooth" structures in carbonates (a bacterial degassing structure) and various types of stromatolites.
Paleogeography and environment of deposition
Paleogeographic reconstructions indicate that the Belt Supergroup accumulated in a fault-bounded rift basin that existed where the North American craton and another landmass were joined in a supercontinent called Columbia/Nuna. The basin appears to have been a closed "lacustrine" environment, or at least not completely open marine. Depositional environments are thought to have ranged from ancient floodplains and exposed mudflats to deep water.
Evidence of the basin-bounding faults exists on all sides of the Belt basin except the west, which rifted away during subsequent continental breakup. The identity of the joined landmass remains controversial. The Siberian craton, Australia and eastern Antarctica have all been suggested based rock ages and paleomagnetic information.
Stratigraphy and distribution
The Belt Supergroup was deposited unconformably on Archean and Paleoproterozoic rocks. It reaches thicknesses of more than and is present in western Montana and northern Idaho, with minor occurrences in northeastern Washington and western Wyoming. Because of this widespread extent, the rock types and formation names of the Belt Supergroup vary depending upon location. In western Montana and northern Idaho the Belt is divided into the following four groups (youngest to oldest):
Missoula Group – Fluvial sands and muds derived from the south.
Piegan Group (Middle Belt Carbonate) – Carbonate muds alternating with laminae of clastic muds.
Ravalli Group – Subaerially deposited sands and muds, mostly fluvial, derived from the southwest.
Lower Belt – Heterogeneous coarse- to fine-grained clastic and carbonate rocks, mostly deep-water deposition with sediments derived from the southwest, and mafic sills.
The Belt Supergroup extends into Canada where the equivalent rocks are called the Purcell Supergroup, and are exposed in southeastern British Columbia and southwestern Alberta.
The Belt Supergroup currently overlies softer Cretaceous age rock that is about 1300 million years younger. This apparent violation of the law of superposition was caused by the Lewis Overthrust.
Economic resources
The Belt Supergroup rocks host a variety of economically significant deposits of lead, zinc, silver, copper, gold, and other metals. These include the Coeur d'Alene lead–zinc–silver mining district in Idaho, which has produced about 7,400,000 tons of lead, 2,900,000 tons of zinc, and 35,600 tons of silver. The equivalent rocks of the Purcell Supergroup in British Columbia include the Sullivan ore body, which has also been a major producer of lead, zinc, and silver.
| Physical sciences | Geologic features | Earth science |
2866340 | https://en.wikipedia.org/wiki/Equipotential | Equipotential | In mathematics and physics, an equipotential or isopotential refers to a region in space where every point is at the same potential. This usually refers to a scalar potential (in that case it is a level set of the potential), although it can also be applied to vector potentials. An equipotential of a scalar potential function in -dimensional space is typically an ()-dimensional space. The del operator illustrates the relationship between a vector field and its associated scalar potential field. An equipotential region might be referred as being 'of equipotential' or simply be called 'an equipotential'.
An equipotential region of a scalar potential in three-dimensional space is often an equipotential surface (or potential isosurface), but it can also be a three-dimensional mathematical solid in space. The gradient of the scalar potential (and hence also its opposite, as in the case of a vector field with an associated potential field) is everywhere perpendicular to the equipotential surface, and zero inside a three-dimensional equipotential region.
Electrical conductors offer an intuitive example. If a and b are any two points within or at the surface of a given conductor, and given there is no flow of charge being exchanged between the two points, then the potential difference is zero between the two points. Thus, an equipotential would contain both points a and b as they have the same potential. Extending this definition, an isopotential is the locus of all points that are of the same potential.
Gravity is perpendicular to the equipotential surfaces of the gravity potential, and in electrostatics and steady electric currents, the electric field (and hence the current, if any) is perpendicular to the equipotential surfaces of the electric potential (voltage).
In gravity, a hollow sphere has a three-dimensional equipotential region inside, with no gravity from the sphere (see shell theorem). In electrostatics, a conductor is a three-dimensional equipotential region. In the case of a hollow conductor (Faraday cage), the equipotential region includes the space inside.
A ball will not be accelerated left or right by the force of gravity if it is resting on a flat, horizontal surface, because it is an equipotential surface.
For the gravity of Earth, the corresponding geopotential isosurface (the equigeopotential) that best fits mean sea level is called the geoid.
| Physical sciences | Classical mechanics | Physics |
2063011 | https://en.wikipedia.org/wiki/Love%20wave | Love wave | In elastodynamics, Love waves, named after Augustus Edward Hough Love, are horizontally polarized surface waves. The Love wave is a result of the interference of many shear waves (S-waves) guided by an elastic layer, which is welded to an elastic half space on one side while bordering a vacuum on the other side. In seismology, Love waves (also known as Q waves (Quer: German for lateral)) are surface seismic waves that cause horizontal shifting of the Earth during an earthquake. Augustus Edward Hough Love predicted the existence of Love waves mathematically in 1911. They form a distinct class, different from other types of seismic waves, such as P-waves and S-waves (both body waves), or Rayleigh waves (another type of surface wave). Love waves travel with a lower velocity than P- or S- waves, but faster than Rayleigh waves. These waves are observed only when there is a low velocity layer overlying a high velocity layer/sub–layers.
Description
The particle motion of a Love wave forms a horizontal line, perpendicular to the direction of propagation (i.e. are transverse waves). Moving deeper into the material, motion can decrease to a "node" and then alternately increase and decrease as one examines deeper layers of particles. The amplitude, or maximum particle motion, often decreases rapidly with depth.
Since Love waves travel on the Earth's surface, the strength (or amplitude) of the waves decrease exponentially with the depth of an earthquake. However, given their confinement to the surface, their amplitude decays only as , where represents the distance the wave has travelled from the earthquake. Surface waves therefore decay more slowly with distance than do body waves, which travel in three dimensions. Large earthquakes may generate Love waves that travel around the Earth several times before dissipating.
Since they decay so slowly, Love waves are the most destructive outside the immediate area of the focus or epicentre of an earthquake. They are what most people feel directly during an earthquake.
In the past, it was often thought that animals like cats and dogs could predict an earthquake before it happened. However, they are simply more sensitive to ground vibrations than humans and are able to detect the subtler body waves that precede Love waves, like the P-waves and the S-waves.
Basic theory
The conservation of linear momentum of a linear elastic material can be written as
where is the displacement vector and is the stiffness tensor. Love waves are a special solution () that satisfy this system of equations. We typically use a Cartesian coordinate system () to describe Love waves.
Consider an isotropic linear elastic medium in which the elastic properties are functions of only the coordinate, i.e., the Lamé parameters and the mass density can be expressed as . Displacements produced by Love waves as a function of time () have the form
These are therefore antiplane shear waves perpendicular to the plane. The function can be expressed as the superposition of harmonic waves with varying wave numbers () and frequencies (). Consider a single harmonic wave, i.e.,
where is the imaginary unit, i.e. . The stresses caused by these displacements are
If we substitute the assumed displacements into the equations for the conservation of momentum, we get a simplified equation
The boundary conditions for a Love wave are that the surface tractions at the free surface must be zero. Another requirement is that the stress component in a layer medium must be continuous at the interfaces of the layers. To convert the second order differential equation in into two first order equations, we express this stress component in the form
to get the first order conservation of momentum equations
The above equations describe an eigenvalue problem whose solution eigenfunctions can be found by a number of numerical methods. Another common, and powerful, approach is the propagator matrix method (also called the matricant approach).
| Physical sciences | Seismology | Earth science |
2063352 | https://en.wikipedia.org/wiki/Mechanised%20agriculture | Mechanised agriculture | Mechanised agriculture or agricultural mechanization is the use of machinery and equipment, ranging from simple and basic hand tools to more sophisticated, motorized equipment and machinery, to perform agricultural operations. In modern times, powered machinery has replaced many farm task formerly carried out by manual labour or by working animals such as oxen, horses and mules.
The entire history of agriculture contains many examples of the use of tools, such as the hoe and the plough. The ongoing integration of machines since the Industrial Revolution has allowed farming to become much less labour-intensive.
Agricultural mechanization is part of this technological evolution of agricultural automation. It can be summarized as a progressive move from manual tools to animal traction, to motorized mechanization, to digital equipment and finally, to robotics with artificial intelligence (AI). These advances can raise productivity and allow for more careful crop, livestock, aquaculture and forestry management; provide better working conditions; improve incomes; reduce the workload of farming; and generate new rural entrepreneurial opportunities.
Current mechanised agriculture includes the use of tractors, trucks, combine harvesters, countless types of farm implements, aeroplanes and helicopters (for aerial application), and other vehicles. Precision agriculture even uses computers in conjunction with satellite imagery and satellite navigation (GPS guidance) to increase yields. New digital equipment is increasingly complementing, or even superseding, motorized machines to make diagnosis and decision-making automatic.
Mechanisation was one of the large factors responsible for urbanisation and industrial economies. Besides improving production efficiency, mechanisation encourages large scale production and sometimes can improve the quality of farm produce. On the other hand, it can cause environmental degradation (such as pollution, deforestation, and soil erosion), especially if it is applied shortsightedly rather than holistically.
History
Jethro Tull's seed drill ( 1701) was a mechanical seed spacing and depth placing device that increased crop yields and saved seed. It was an important factor in the British Agricultural Revolution.
Since the beginning of agriculture threshing was done by hand with a flail, requiring a great deal of labour. The threshing machine, which was invented in 1794 but not widely used for several more decades, simplified the operation and allowed the use of animal power. Before the invention of the grain cradle (ca. 1790) an able bodied labourer could reap about one quarter acre of wheat in a day using a sickle. It was estimated that each of Cyrus McCormick's horse-pulled reapers (ca. 1830s) freed up five men for military service in the US Civil War. Later innovations included raking and binding machines. By 1890 two men and two horses could cut, rake and bind 20 acres of wheat per day.
In the 1880s the reaper and threshing machine were combined into the combine harvester. These machines required large teams of horses or mules to pull. Steam power was applied to threshing machines in the late 19th century. There were steam engines that moved around on wheels under their own power for supplying temporary power to stationary threshing machines. These were called road engines, and Henry Ford seeing one as a boy was inspired to build an automobile.
With internal combustion came the first modern tractors in the early 1900s, becoming more popular after the Fordson tractor (ca. 1917). At first reapers and combine harvesters were pulled by teams of horses or tractors, but in the 1930s self powered combines were developed.
Advertising for motorised equipment in farm journals during this era did its best to compete against horse-drawn methods with economic arguments, extolling common themes such as that a tractor "eats only when it works", that one tractor could replace many horses, and that mechanisation could allow one man to get more work done per day than he ever had before. The horse population in the US began to decline in the 1920s after the conversion of agriculture and transportation to internal combustion. Peak tractor sales in the US were around 1950. In addition to saving labour, this freed up much land previously used for supporting draft animals. The greatest period of growth in agricultural productivity in the US was from the 1940s to the 1970s, during which time agriculture was benefiting from internal combustion powered tractors and combine harvesters, chemical fertilisers and the green revolution.
Although U.S. farmers of corn, wheat, soy, and other commodity crops had replaced most of their workers with harvesting machines and combines by the 1950s enabling them to efficiently cut and gather grains, growers of produce continued to rely on human pickers to avoid the bruising of the product in order to maintain the blemish-free appearance demanded by customers. The continuous supply of undocumented workers from Latin America that harvest the crops for low wages further suppressed the need for mechanisation. As the number of undocumented workers has continued to decline since reaching its peak in 2007 due to increased border patrols and an improving Mexican economy, the industry is increasing the use of mechanisation. Proponents argue that mechanisation will boost productivity and help to maintain low food prices while farm worker advocates assert that it will eliminate jobs and will give an advantage to large growers who are able to afford the required equipment.
Motorized mechanization adoption trends
Motorized mechanization has substantially expanded at global level, although it has been unevenly and inadequately adopted particularly in sub-Saharan Africa. Mechanization is limited to a range of operations including harvesting and weeding and is rarely used for fruit and vegetable production across the globe.
Extensive adoption started in the United States of America, where tractors replaced about 24 million draught animals between 1910 and 1960 and become the main source of farm power. United Kingdom first started using tractors in the 1930s, but agricultural transformation in Japan and some European countries (Denmark, France, Germany, Spain and former Yugoslavia) did not take place until about 1955. Thereafter, the adoption of motorized mechanization took place very quickly, completely superseding animal traction. Using tractors as farm power enabled, and even triggered, innovations in other agricultural machinery and equipment that greatly eased the toil associated with agriculture and allowed farmers to carry out tasks more quickly. At a later stage, motorized machinery also increased in many Asian and Latin American countries.
Sub-Saharan Africa is the only region where adoption of motorized mechanization has not progressed over the past decades. A study in 11 countries proves this low level of mechanization in the region, finding that only 18 percent of the sampled households have access to tractor-powered appliances. The remaining ones make use of either simple hand-held tools (48 percent) or animal-powered equipment (33 percent).
Employment impact
Since at least the early nineteenth century there have been concerns over the possible negative socioeconomic impacts of labour-saving technological change, particularly job displacement resulting in unemployment. However, fears that automation increases labour productivity to the extent that it causes massive unemployment are not supported by historical realities. Instead, innovation and incorporation of labour-saving technologies tends to take long, and automation of one task often spurs increases in the need for workers to perform other jobs. The direct impact of automation on employment will be determined by the factors leading to its adoption.
If rising wages and labour scarcities drive the adoption of automation then it is not likely to create unemployment. Automation can also stimulate agricultural employment. For example, it can enable farms to increase their production following growing food demand. Agricultural automation is part of the structural transformation of societies through which increased agricultural labour productivity gradually releases agricultural workers, giving them the opportunity to take new jobs in other sectors, including industry and services. On the other hand, automation that is forcibly promoted, such as through government subsidies, could cause rising unemployment and falling or stagnant wages.
The Food and Agriculture Organization of the United Nations (FAO) advises against governments implementing distortive subsidies for automation because doing so risks increasing unemployment. FAO also advises against restricting automation on the assumption that this will save jobs and incomes, because it risks making agriculture less competitive and productive. Instead, the recommendation is to concentrate on creating an enabling environment to adopt automation – particularly by small-scale agricultural producers, women and youth – while making social protection available to least skilled workers, who are more likely to lose their jobs during the transition.
Applications
Preparing land for planting
Seed drilling, planting
It is done by the seed drill. The plantation of seeds depends upon the season.
Weeding, crop spraying
Harvesting
Asparagus are presently harvested by hand with labor costs at 71% of production costs and 44% of selling costs. Asparagus is a difficult crop to harvest since each spear matures at a different speed making it difficult to achieve a uniform harvest. A prototype asparagus harvesting machineusing a light-beam sensor to identify the taller spearsis expected to be available for commercial use.
Mechanization of Maine's blueberry industry has reduced the number of migrant workers required from 5,000 in 2005 to 1,500 in 2015 even though production has increased from 50 to 60 million pounds per year in 2005 to 90 million pounds in 2015.
As of 2014, prototype chili pepper harvesters are being tested by New Mexico State University. The New Mexico green chile crop is currently hand-picked entirely by field workers as chili pods tend to bruise easily. The first commercial application commenced in 2015. The equipment is expected to increase yield per acre and help to offset a sharp decline in acreage planted due to the lack of available labour and drought conditions.
As of 2010, approximately 10% of the processing orange acreage in Florida is harvested mechanically, mainly with citrus canopy shaker machines. Mechanization has progressed slowly due to the uncertainty of future economic benefits due to competition from Brazil and the transitory damage to orange trees when they are harvested.
There has been an ongoing transition to mechanical harvesting of cling peaches (mostly used in canning) where the cost of labor is 70 percent of a grower's direct costs. In 2016, 12 percent of the cling peach tonnage from Yuba County and Sutter County in California will be mechanically harvested. Fresh peaches destined for direct customer sales must still be hand-picked.
As of 2007, mechanized harvesting of raisins is at 45%; however the rate has slowed due to high raisin demand and prices making the conversion away from hand labour less urgent. A new strain of grape developed by the USDA that drys on the vine and is easily harvested mechanically is expected to reduce the demand for labour.
Strawberries are a high cost-high value crop with the economics supporting mechanization. In 2005, picking and hauling costs were estimated at $594 per ton or 51% of the total grower cost. However, the delicate nature of fruit make it an unlikely candidate for mechanization in the near future. A strawberry harvester developed by Shibuya Seiki and unveiled in Japan in 2013 is able to pick a strawberry every eight seconds. The robot identifies which strawberries are ready to pick by using three separate cameras and then once identified as ready, a mechanized arm snips the fruit free and gently places it in a basket. The robot moves on rails between the rows of strawberries which are generally contained within elevated greenhouses. The machine costs 5 million yen. A new strawberry harvester made by Agrobot that will harvest strawberries on raised, hydroponic beds using 60 robotic arms is expected to be released in 2016.
Mechanical harvesting of tomatoes started in 1965 and as of 2010, nearly all processing tomatoes are mechanically harvested. As of 2010, 95% of the US processed tomato crop is produced in California. Although fresh market tomatoes have substantial hand harvesting costs (in 2007, the costs of hand picking and hauling were $86 per ton which is 19% of total grower cost), packing and selling costs were more of a concern (at 44% of total grower cost) making it likely that cost saving efforts would be applied there.
According to a 1977 report by the California Agrarian Action Project, during the summer of 1976 in California, many harvest machines had been equipped with a photo-electric scanner that sorted out green tomatoes among the ripe red ones using infrared lights and colour sensors. It worked in lieu of 5,000 hand harvesters causing displacement of innumerable farm labourers as well as wage cuts and shorter work periods. Migrant workers were hit the hardest. To withstand the rigour of the machines, new crop varieties were bred to match the automated pickers. UC Davis Professor G.C. Hanna propagated a thick-skinned tomato called VF-145. But even still, millions were damaged with impact cracks and university breeders produced a tougher and juiceless "square round" tomato. Small farms were of insufficient size to obtain financing to purchase the equipment and within 10 years, 85% of the state's 4,000 cannery tomato farmers were out of the business. This led to a concentrated tomato industry in California that "now packed 85% of the nation’s tomato products". The monoculture fields fostered rapid pest growth, requiring the use of "more than four million pounds of pesticides each year" which greatly affected the health of the soil, the farm workers, and possibly the customers.
| Technology | Agriculture_2 | null |
2063359 | https://en.wikipedia.org/wiki/Chi%20%28unit%29 | Chi (unit) | The chi (Tongyong Pinyin chih) is a traditional Chinese unit of length. Although it is often translated as the "", its length was originally derived from the distance measured by a human hand, from the tip of the thumb to the tip of the forefinger, and is similar to the ancient span. It first appeared during China's Shang dynasty approximately 3,000 years ago and has since been adopted by other East Asian cultures such as Japan (shaku), Korea (ja/cheok), and Vietnam (thước). Its present value is standardized at around , although the exact standards vary among the mainland of the People's Republic of China, its special administrative region of Hong Kong, and Taiwan.
In its ancient and modern forms, the chi is divided into 10 smaller units known as cun (the "Chinese inch"). 10 chi are equal to 1 zhàng.
Modern values
In the People's Republic of China, since 1984, the chi has been defined as exactly 1/3 of a metre, i.e., . However, in the Hong Kong SAR the corresponding unit, pronounced tsek (cek3) in Cantonese, is defined as exactly or 1 7/32 ft. The two units are sometimes referred to in English as "Chinese foot" and "Hong Kong foot".
In Taiwan, chi is the same as the Japanese shaku, i.e., .
Historical values
The study of ancient rulers and other artifacts whose size in the contemporary chi was known allowed modern researchers to surmise that during the 2nd century BC to 3rd century AD the (Qin dynasty to Han dynasty to the Three Kingdoms period), the value of the chi varied between . Even earlier, during the Warring States era, the value of chi was essentially the same.
It is thought that the ancient Chinese astronomers also used chi as an angular unit; modern analysis of historical records indicates that it may have been equal to one degree.
In the 19th century, the value of the chi, depending on the part of the country and the application, varied between . According to an 1864 British report, in most of China the chi used by engineers in public works was equal to , the surveyors' chi was , while the value generally used for measuring distances was . In Guangzhou, however, the chi used for local trade varied from – i.e., very close to the modern chek. The value fixed by a Sino-British treaty for the purposes of customs duties in Hong Kong was .
In 1905 the Korean Empire defined the cheok as 10000/33000 of a metre. In 1964, South Korea fully adopted SI units.
Usage in Chinese language
Due to its long history and its widespread usage, chi (along with cun) has also seen metaphorical usages in the Chinese language. For example, chi cun (), a word made up of the units chi and cun, refers to the dimensions of an object, while the idiom "dé cùn jìn chǐ" () means "extremely greedy".
In informal use in China, chi is also sometimes used to refer to the United States customary foot or British imperial foot.
| Physical sciences | East Asian | Basics and measurement |
2064083 | https://en.wikipedia.org/wiki/Layering | Layering | Layering is a vegetative propagation technique where the stem or branch of a plant is manipulated to promote root development while still attached to the parent plant. Once roots are established, the new plant can be detached from the parent and planted. Layering is utilized by horticulturists to propagate desirable plants.
Natural layering typically occurs when a branch touches the ground, whereupon it produces adventitious roots. At a later stage the connection with the parent plant is severed and a new plant is produced as a result.
The horticultural layering process typically involves wounding the target region to expose the inner stem and optionally applying rooting compounds.
In ground layering or simple layering, the stem is bent down and the target region is buried in the soil. This is done in plant nurseries in imitation of natural layering by many plants such as brambles which bow over and touch the tip on the ground, at which point it grows roots and, when separated, can continue as a separate plant. In either case, the rooting process may take from several weeks to a year.
There are two methods of air layering, which do not involve burying the stem. In ring air layering, the exposed wound is covered in a growth medium such as sphagnum moss, and wrapped in a material such as plastic. The roots grow into the medium and after a period of time, the stem is separated from the original plant. Tourniquet air layering has a similar method to air layering, except that instead of creating a wound, a wire is wrapped around the stem and the ends are twisted until it is very tight.
Layering is more complicated than taking cuttings, but has the advantage that the propagated portion continues to receive water and nutrients from the parent plant while it is forming roots. This is important for plants that form roots slowly, or for propagating large pieces. Layering is used quite frequently in the propagation of bonsai; it is also used as a technique for both creating new roots and improving existing roots.
Method and process
A low-growing stem is bent down to touch a hole dug in the ground, then pinned in place using something shaped like a clothes hanger hook and covered over with soil. However, a few inches of leafy growth must remain above the ground for the bent stem to grow into a new plant. Removing a section of skin from the lower-facing stem part before burying may help the rooting process. If using rooting hormone, the stem should be cut just beneath a node. The resultant notch should be wedged open with a toothpick or similar piece of wood and the hormone applied before burying.
The buried stem part then grows new roots which take hold in the soil while the above parts grow as individual new plants. Once the end of the stem has grown long enough the process can be repeated, creating the appearance of a row of plants linked by humped, intermittently buried stems. Better results can be achieved when the top of the plant is closer to the vertical.
Once the process is completed, the buried section should be kept well-watered until roots begin to form. The new individual plant may require one to two years before it is strong enough to survive on its own. When it is, the original stem should be cut where it enters the ground, thereby separating the two plants.
Plants types and benefit
As layering does not involve sexual reproduction, new plants are effectively clones of the original plant and will exhibit the same characteristics. This includes flower, fruit and foliage. Plant selection usually involves plants with a flexible stem.
Simple layering can be more attractive when managing a cascading or spreading plant. These plants tend to propagate in this manner anyway, and potting a new limb will give extra plants without having to sow new seed.
Simple layering can also help when a plant has overgrown its pot and is drooping over the side. The long stem is layered into another pot until it roots, thus bringing it back to soil level.
Ground layering
Ground layering or mound layering is the typical propagation technique for the popular Malling-Merton series of clonal apple root stocks, in which the original plants are set in the ground with the stem nearly horizontal, which forces side buds to grow upward. After these are started at the start of the growing season, the original stem is buried up to some distance from the tip. At the end of the growing season, the side branches will have rooted, and can be separated while the plant is dormant. Some of these will be used for grafting rootstocks, and some can be reused in the nursery for the next growing season's crop.
Ground layering is used in the formation of visible surface roots, known as "nebari", on bonsai trees.
Air layering
In air layering (or marcotting), the target region is wounded by an upward 4 cm long cut and held open with a toothpick or similar, or a strip of bark is removed. The wound is then surrounded with a lump of moisture-retaining medium such as sphagnum moss or cloth, and then further surrounded by a moisture barrier such as plastic film tied or taped to the branch to prevent moisture loss or ingress of too much water as from rain. Rooting hormone is often applied to the wound to encourage root growth. When sufficient roots have grown from the wound, the stem is removed from the parent plant and planted, taking care to shield it from too much sun and to protect it from drying out until the new roots take hold. It can take the layer from a few weeks to one or more growing seasons to produce sufficient roots; this is largely dependent on the plant species and the vigor of the parent plant.
| Technology | Horticulture | null |
2065886 | https://en.wikipedia.org/wiki/BET%20theory | BET theory | Brunauer–Emmett–Teller (BET) theory aims to explain the physical adsorption of gas molecules on a solid surface and serves as the basis for an important analysis technique for the measurement of the specific surface area of materials. The observations are very often referred to as physical adsorption or physisorption. In 1938, Stephen Brunauer, Paul Hugh Emmett, and Edward Teller presented their theory in the Journal of the American Chemical Society. BET theory applies to systems of multilayer adsorption that usually utilizes a probing gas (called the adsorbate) that does not react chemically with the adsorptive (the material upon which the gas attaches to) to quantify specific surface area. Nitrogen is the most commonly employed gaseous adsorbate for probing surface(s). For this reason, standard BET analysis is most often conducted at the boiling temperature of N2 (77 K). Other probing adsorbates are also utilized, albeit less often, allowing the measurement of surface area at different temperatures and measurement scales. These include argon, carbon dioxide, and water. Specific surface area is a scale-dependent property, with no single true value of specific surface area definable, and thus quantities of specific surface area determined through BET theory may depend on the adsorbate molecule utilized and its adsorption cross section.
Concept
The concept of the theory is an extension of the Langmuir theory, which is a theory for monolayer molecular adsorption, to multilayer adsorption with the following hypotheses:
gas molecules physically adsorb on a solid in layers infinitely;
gas molecules only interact with adjacent layers; and
the Langmuir theory can be applied to each layer.
the enthalpy of adsorption for the first layer is constant and greater than the second (and higher).
the enthalpy of adsorption for the second (and higher) layers is the same as the enthalpy of liquefaction.
The resulting BET equation is
where c is referred to as the BET C-constant, is the vapor pressure of the adsorptive bulk liquid phase which would be at the temperature of the adsorbate and θ is the surface coverage, defined as:
.
Here is the amount of adsorbate and is called the monolayer equivalent. The is the entire amount that would be present as a monolayer (which is theoretically impossible for physical adsorption) that would cover the surface with exactly one layer of adsorbate. The above equation is usually rearranged to yield the following equation for the ease of analysis:
where and are the equilibrium and the saturation pressure of adsorbates at the temperature of adsorption, respectively; is the adsorbed gas quantity (for example, in volume units) while is the monolayer adsorbed gas quantity. is the BET constant,
where is the heat of adsorption for the first layer, and is that for the second and higher layers and is equal to the heat of liquefaction or heat of vaporization.
Equation (1) is an adsorption isotherm and can be plotted as a straight line with on the y-axis and on the x-axis according to experimental results. This plot is called a BET plot. The linear relationship of this equation is maintained only in the range of . The value of the slope and the y-intercept of the line are used to calculate the monolayer adsorbed gas quantity and the BET constant . The following equations can be used:
The BET method is widely used in materials science for the calculation of surface areas of solids by physical adsorption of gas molecules. The total surface area and the specific surface area are given by
where is in units of volume which are also the units of the monolayer volume of the adsorbate gas,
is the Avogadro number, the adsorption cross section of the adsorbate, the molar volume of the adsorbate gas, and the mass of the solid sample or adsorbent.
Derivation
The BET theory can be derived similarly to the Langmuir theory, but by considering multilayered gas molecule adsorption, where it is not required for a layer to be completed before an upper layer formation starts. Furthermore, the authors made five assumptions:
Adsorptions occur only on well-defined sites of the sample surface (one per molecule)
The only molecular interaction considered is the following one: a molecule can act as a single adsorption site for a molecule of the upper layer.
The uppermost molecule layer is in equilibrium with the gas phase, i.e. similar molecule adsorption and desorption rates.
The desorption is a kinetically limited process, i.e. a heat of adsorption must be provided:
these phenomena are homogeneous, i.e. same heat of adsorption for a given molecule layer.
it is E1 for the first layer, i.e. the heat of adsorption at the solid sample surface
the other layers are assumed similar and can be represented as condensed species, i.e. liquid state. Hence, the heat of adsorption is EL is equal to the heat of liquefaction.
At the saturation pressure, the molecule layer number tends to infinity (i.e. equivalent to the sample being surrounded by a liquid phase)
Consider a given amount of solid sample in a controlled atmosphere. Let θi be the fractional coverage of the sample surface covered by a number i of successive molecule layers. Let us assume that the adsorption rate Rads,i-1 for molecules on a layer (i-1) (i.e. formation of a layer i) is proportional to both its fractional surface θi-1 and to the pressure P, and that the desorption rate Rdes,i on a layer i is also proportional to its fractional surface θi:
where ki and k−i are the kinetic constants (depending on the temperature) for the adsorption on the layer (i−1) and desorption on layer i, respectively. For the adsorptions, these constants are assumed similar whatever the surface.
Assuming an Arrhenius law for desorption, the related constants can be expressed as
where Ei is the heat of adsorption, equal to E1 at the sample surface and to EL otherwise.
Consider some substance A. The adsorption of A onto an available surface site produces a new site on the first layer. In summary,
A(g) + <=>
Extending this to higher order layers one obtains
A(g) + <=>
and similarly
A(g) + <=>
Denoting the activity of the number of available sites of the th layer with and the partial pressure of A with , the last equilibrium can be written
It follows that the coverage of the first layer can be written
and that the coverage of the second layer can be written
Realising that the adsorption of A onto the second layer is equivalent to adsorption of A onto its own liquid phase, the rate constant for should be the same, which results in the recursion
In order to simplify some infinite summations, let and let . Then the th layer coverage can written
if . The coverage of any layer is defined as the relative number of available sites. An alternative definition, which leads to a set of coverage's that are numerically to those resulting from the original way of defining surface coverage, is that denotes the relative number of sites covered by only adsorbents. Doing so it is easy to see that the total volume of adsorbed molecules can be written as the sum
where is the molecular volume. Employing the fact that this sum is the first derivative of a geometric sum, the volume becomes
Since the total coverage of a mono-layer must be unity, the full mono-layer coverage must be
In order to properly make the substitution for , the restriction forces us to take the zeroth contribution outside the summation, resulting in
Lastly, defining the excess coverage as , the excess volume relative to the volume of an adsorbed mono-layer becomes
where the last equality was obtained by making use of the series expansions presented above. The constant must be interpreted as the relative binding affinity the substance A has towards a surface, relative to its own liquid. If then the initial part of the isotherm will be reminiscent of the Langmuir isotherm which reaches a plateau at full mono-layer coverage, whereas means the mono-layer will have a slow build-up. Another thing to note is that in order for the geometric substitutions to hold, . The isotherm above exhibits a singularity at . Since one can write , implying that . This means that must be true, ultimately resulting in .
Finding the linear BET range
It is still not clear on how to find the linear range of the BET plot for microporous materials in a way that reduces any subjectivity in the assessment of the monolayer capacity. A crowd-sourced study involving 61 research groups has shown that reproducibility of BET area determination from identical isotherms is, in some cases, problematic. Rouquerol et al. suggested a procedure that is based on two criteria:
C must be positive implying that any negative intercept on the BET plot indicates that one is outside the valid range of the BET equation.
Application of the BET equation must be limited to the range where the term V(1-P/P0) continuously increases with P/P0.
These corrections are an attempt to salvage the BET theory, which is restricted to type II isotherms. Even while using this type, use of the data itself is restricted to 0.05 to 0.35 of , routinely discarding 70% of the data. This restriction must be modified depending upon conditions.
Limitations of BET
Terrell L. Hill described BET as a theory that is "... extremely useful as a qualitative guide; but it is not quantitatively correct". Although BET adsorption isotherm is still extensively used for different applications and is used for specific surface area determinations of powders whose calculation is not sensitive to the simplifications of the BET theory. Both Hackerman's and Sing's group have highlighted the limitations of the BET method. Hackerman et al. noted the potential for 10% uncertainty in the method's values, with Sing's group attributed the significant variation in reported values of molecular area to the BET method's possible inaccurate assessment of monolayer capacity. In subsequent studies using the BET interpretation of nitrogen and water vapor adsorption isotherms, the reported area occupied by an adsorbed water molecule on fully hydroxylated silica ranged from 0.25 to 0.44 nm². Other issues with the BET include the fact that in certain cases, BET leads to anomalies such as reaching an infinite amount adsorbed at reaching unity, and in some cases, the constant C (surface binding energy) can be determined to be negative.
Applications
Cement and concrete
The rate of curing of concrete depends on the fineness of the cement and of the components used in its manufacture, which may include fly ash, silica fume and other materials, in addition to the calcinated limestone which causes it to harden. Although the Blaine air permeability method is often preferred, due to its simplicity and low cost, the nitrogen BET method is also used.
When hydrated cement hardens, the calcium silicate hydrate (or C-S-H), which is responsible for the hardening reaction, has a large specific surface area because of its high porosity. This porosity is related to a number of important properties of the material, including the strength and permeability, which in turn affect the properties of the resulting concrete. Measurement of the specific surface area using the BET method is useful for comparing different cements. This may be performed using adsorption isotherms measured in different ways, including the adsorption of water vapour at temperatures near ambient, and adsorption of nitrogen at 77 K (the boiling point of liquid nitrogen). Different methods of measuring cement paste surface areas often give very different values, but for a single method the results are still useful for comparing different cements.
Activated carbon
Activated carbon has strong affinity for many gases and has an adsorption cross section of 0.162 nm2 for nitrogen adsorption at liquid-nitrogen temperature (77 K). BET theory can be applied to estimate the specific surface area of activated carbon from experimental data, demonstrating a large specific surface area, even around 3000 m2/g. However, this surface area is largely overestimated due to enhanced adsorption in micropores, and more realistic methods should be used for its estimation, such as the subtracting pore effect (SPE) method.
Catalysis
In the field of solid catalysis, the surface area of catalysts is an important factor in catalytic activity. Inorganic materials such as mesoporous silica and layered clay minerals have high surface areas of several hundred m2/g calculated by the BET method, indicating the possibility of application for efficient catalytic materials.
Specific surface area calculation
The ISO 9277 standard for calculating the specific surface area of solids is based on the BET method. The method has also been adapted for determination of specific area of ceramics and non-ferrous metal powders.
Thermal desorption
In 2023, researchers in the United States developed a method to determine BET surface areas using a thermogravimetric analyzer (TGA). This method uses a TGA to heat a porous sample loaded with an adsorbate, the produced plot of sample mass vs. temperature is then mapped into a standard isotherm to which BET theory is applied as normal. Common fluids, e.g. water or toluene, can be used as adsorbates for the TGA method allowing the specific interactions of different adsorbates to be determined, as these frequently differ from the commonly used nitrogen.
| Physical sciences | Other separations | Chemistry |
2067207 | https://en.wikipedia.org/wiki/Kolkata%20Metro | Kolkata Metro | The Kolkata Metro is a rapid transit system serving the city of Kolkata and the Kolkata Metropolitan Region in West Bengal, India. Opened in 1984, it was the first operational rapid transit system in India. Currently it is the second busiest and fifth longest rapid transit system in India. As of now, it has one fully operational, and three partly operational lines for a total of . Two other lines are in various phases of construction and planning. The system has a mix of underground, at-grade, and elevated stations using both broad-gauge and standard-gauge tracks. It operates on a 750 V DC Third rail system. Trains operate between 06:50 and 22:40 IST and the fares range from ₹5 to ₹50.
The Kolkata Metro was initially planned in the 1920s, but construction started in the 1970s. The first underground stretch, from Bhawanipore (now Netaji Bhawan) to Esplanade, opened in 1984. A truncated section of Line 2, or the East–West Corridor, from Salt Lake Sector V to Phoolbagan opened in 2020. Line 3, or the Joka-Esplanade Corridor (currently truncated in Majerhat), opened in 2022 while Line 6, from Kavi Subhash to Hemanta Mukhopadhyay, opened in 2024. It is the second busiest metro network in India after the Delhi Metro and is the fifth-longest operational metro network in India after the Delhi Metro, Namma Metro, Hyderabad Metro, and Mumbai Metro.
Metro Railway, Kolkata and Kolkata Metro Rail Corporation are the owners and operator of the system. On 29 December 2010, Metro Railway, Kolkata, became the 17th zone of the Indian Railways, completely owned and funded by the Ministry of Railways. It is the only metro system in the country to be controlled by Indian Railways. Around 300 daily train trips carry more than passengers.
History
Early attempts
In the September 1919 session of the Imperial Legislative Council at Shimla, a committee was set up by W. E. Crum that recommended a metro line for Kolkata (formerly Calcutta). This line was supposed to connect Bagmari in the east to Benaras Road, Salkia, in Howrah in the west via a tunnel beneath Hooghly River. The estimated construction costs were £3,526,154, about based on current exchange rates, and the proposed deadline was 1925–1926. The proposed line was long, about shorter than the current East-West Corridor, which would connect East Bengal Railway in Bagmari and East Indian Railway in Benaras Road. The tickets were priced at 3 annas ( 0.1875) for the full trip. Crum also mentioned a north–south corridor back then. An east–west metro railway connection, named the "East–West Tube Railway", was proposed for Kolkata in 1921 by Harley Dalrymple-Hay. All the reports can be found in his 1921 book Calcutta Tube Railways. However, in 1923, the proposal was not undertaken due to a lack of funds.
Planning
Then the Chief Minister of West Bengal, Bidhan Chandra Roy, reconceived the idea of an underground railway for Kolkata from 1949 to 1950. A team of French experts conducted a survey, but nothing concrete materialized. Efforts to solve the traffic problem by augmenting the existing fleet of public transport vehicles hardly helped, since roads accounted for only 4.2 percent of the surface area in Kolkata, compared with 25 percent in Delhi and 30 percent in other cities. To find alternative solutions, the Metropolitan Transport Project (MTP) was set up in 1969. The MTP, with the help of Soviet specialists, Lenmetroproekt and East German engineers, prepared a master plan to provide five rapid-transit (metro) lines for the city of Kolkata, totaling a length of , in 1971. Three were selected for construction. These were:
Dum Dum – Tollygunge (Line 1. Presently operates from Dakshineswar to New Garia)
Bidhannagar – Ramrajatala (Line 2. Presently truncated till Howrah Maidan)
Dakshineswar – Thakurpukur (Divided into Line 1; Noapara to Dakshineswar and Line 3; Joka to Esplanade)
The highest priority was given to the busy north–south corridor between Dum Dum and Tollygunge over a length of ; work on this project was approved on 1 June 1972. A tentative deadline was fixed to complete all the corridors by 1991.
Construction
North-South Metro
Since it was India's first metro and was constructed as a completely indigenous process, a traditional cut-and-cover method and driven shield tunneling was chosen and the Kolkata Metro was more of a trial-and-error affair, in contrast to the Delhi Metro, which saw the involvement of multiple international consultants. As a result, it took nearly 23 years to completely construct the underground railway.
The foundation stone of the project was laid by Indira Gandhi, the Prime Minister of India, on 29 December 1972, and construction work started in 1973–74. Initially, cut and cover along with slurry wall construction to handle soft ground, was recommended by the Soviet Union consultants. Later, in 1977, it was decided to adopt both shield tunneling and cut and cover methods for the construction of underpopulated areas, sewer lines, water mains, electrical cables, telephone cables, tram lines, canals, etc. The technology was provided by M/s NIKEX Hungarian Co., Budapest. In the early days, the project was led by the Union Railway Minister from West Bengal, A. B. A. Ghani Khan Choudhury, often against the prevailing socio-political stance of his contemporaries in the West Bengal government. From the start of construction, the project had to contend with several problems including insufficient funds (until 1977–1978), a shifting of underground utilities, court injunctions, and an irregular supply of vital materials. In 1977, an injunction for the allocation of new funding was passed by the newly elected Jyoti Basu government.
Despite all the hurdles, services began on 24 October 1984, with the commissioning of a partial commercial service covering a distance of with five stations served between Esplanade and Bhowanipur (currently Netaji Bhavan). The first metro was driven by Tapan Kumar Nath and Sanjoy Kumar Sil. The service was quickly followed by commuter services on another stretch in the north between Dum Dum and Belgachhia on 12 November 1984. The commuter service was extended to Tollygunge on 29 April 1986, covering a further distance of , making the service available over a distance of and covering 11 stations. However, the services on the north section were suspended starting 26 October 1992, as this small, isolated section was little used. The Line 1 was almost entirely built by cut and cover method, while a small 1.09 km stretch between Belgachia and Shyambazar was built using shield tunneling with compressed air and air locks, since the alignment crossed a railway yard (now Kolkata railway station) and Circular Canal.
After more than eight years, the Belgachhia–Shyambazaar section, along with the Dum Dum–Belgachhia stretch, was opened on 13 August 1994. Another stretch from Esplanade to Chandni Chowk was commissioned shortly afterward, on 2 October 1994. The Shyambazaar-Shobhabazar–Girish Park () and Chandni Chowk–Central () sections were opened on 19 February 1995. Services on the entire stretch of the Metro were introduced from 27 September 1995 by bridging the gap with Mahatma Gandhi Road metro station in the middle.
In 1999–2000, the extension of Line 1 along an elevated corridor from Tollygunge to New Garia, with six stations, was sanctioned at a cost of . The section was constructed and opened in two phases, Mahanayak Uttam Kumar to Kavi Nazrul in 2009 and Kavi Nazrul to Kavi Subhash in 2010. In the north, the line was extended till Noapara from Dum Dum on 10 July 2013. The latest extension opened was the stretch from Noapara to Dakshineswar on 23 February 2021.
East-West Metro
The master plan of the metro corridor was made in 1971 along with the North–South Corridor, connecting the office district of Bidhannagar with the twin city and transportation hub Howrah via another transport hub of the city, Sealdah, and the central business district Esplanade by an underwater metro line. It is a project, sanctioned in 2008 by Prime Minister Manmohan Singh. The foundation stone was laid on 22 February 2009 and construction started in March 2009. The autonomous Kolkata Metro Rail Corporation (KMRC) was formed to implement the project. The Government of India (Ministry of Urban Development) and Government of West Bengal each had a half-share in it. Later, the Government of West Bengal pulled out from it, and the shares were transferred to the Ministry of Railways.
Route realignment
The realignment led to many other issues and delays. Some of the biggest issues were the H-piles under the Esplanade metro station and the Bowbazar mishap. Per the 1971 master plan, the East-West Corridor was supposed to pass under Central metro station, so the square foundational beams in Esplanade were not removed. Since the Tunnel Boring Machines (TBMs) cannot cut through steel, another small tunnel was dug using New Austrian tunneling method (NATM) and the H-piles were cut manually. This extended the tunneling process by one and a half months. In September 2019, during the construction of the eastbound tunnel (from Esplanade to Sealdah), a TBM hit an aquifer under Bowbazar, causing a major collapse in the area, delaying work in that section for several months. Around 80 houses were damaged and many buildings were declared unsafe, affecting more than 600 people. Later subsidence in the area was checked using grouting.
Expansion planning
By 2011–2012, the Railway Ministry had announced plans for the construction of five new metro lines and an extension of the existing north–south corridor. These were:
Salt Lake – Howrah Maidan (Line 2 or East–West Metro Corridor)
Joka – B.B.D. Bagh (Line 3. Later truncated till Esplanade)
Noapara – Barasat (Line 4, via airport)
Baranagar – Barrackpore (Line 5)
New Garia – Dum Dum Airport (Line 6)
Major modifications
A new four-platform interchange station was constructed at Noapara and Kavi Subhash. This acts as an interchange station for Line 1 with Line 4 and Line 6 respectively. For the time being, only two platforms are in use for Noapara, but once Line 4 is running, all four platforms will be operational, whereas, from 6 March 2024, all four platforms of Kavi Subhash will be operational.
The existing Esplanade metro station was upgraded and a subway was constructed to the new metro station to provide an interchange among Line 1, Line 2 and future Line 3. In 2009–2010, Line 1 underwent upgrades of services and amenities and many stations were renamed after famous personalities by then Minister of Railways Mamata Banerjee.
Network
Summary
The Kolkata Metro currently operates with four lines: Line 1, Line 2, Line 3, and Line 6. These lines have a total of 50 operational stations with a further 29 under construction.
Line 1, also known as the north–south line connects Kavi Subhash to Dakshineswar and consists of 26 stations.
Line 2, also known as the East-West Line, currently connects Salt Lake Sector V to Sealdah and Esplanade with Howrah Maidan, with a gap between Sealdah and Esplanade. This line currently consists of 12 stations and proposals have been made to extend it from Sector-V to Teghoria.
Line 3 presently connects Joka to Majerhat, with 7 stations in operation and construction is ongoing up to Esplanade. The line is also planned to extend from Joka to Diamond Park.
Line 4 is currently under trials from Noapara to Biman Bandar(Kolkata Airport) and under construction from Biman Bandar to New Barrackpore with a planned extension till Barasat .
Line 6 currently connects Kavi Subhash to Hemanta Mukhopadhyay consisting of 5 stations. This line is under construction up to Jai Hind (Kolkata Airport) and when finished will have a total of 24 stations.
Additionally, there is one more line proposed,
Line 5 which will connect Baranagar to Barrackpore with a planned 11 stations.
The planned network will have a total of 21 stations, further expanding the metro's reach and capacity. The planned network consists of the entirety of Line 5; the extension of Line 2 to Teghoria and the extension of Line 4 to Barasat.
Operational
Under construction/planned
Lines
Summary
Line 1, or the Blue line, of Kolkata Metro () has a total length of serving 26 Kolkata Metro stations, of which 15 are underground, 9 are elevated and 2 at-grade. It uses the broad gauge tracks. It was the first underground railway to be built in India, with the first trains running in October 1984 and the full stretch that had been initially planned completed and operational by February 1995. The southward extension of the Blue Line to an elevated corridor from Tollygunge to New Garia was constructed and opened in two phases, Mahanayak Uttam Kumar to Kavi Nazrul in 2009 and Kavi Nazrul to Kavi Subhash in 2010. Another extension constructed was the elevated corridor from Dum Dum to Noapara in 2013. The last extension from Noapara to Dakshineswar opened in 2021, thus completing the Blue line.
A northward extension from Dum Dum to Dakshineswar () was sanctioned and included in the 2010–2011 budget at a cost of . The commercial operations for Dum Dum to Noapara () were commissioned in March 2013, and construction from Noapara to Dakshineswar with an interchange with Line 5 at Baranagar () is being executed by RVNL. This section is opened on 23 February 2021 for general public with a projected ridership of 55,000 by 2030.
An upgrade of the existing signaling system from Indian Railways Signalling to Communication Based Train Control was proposed by Metro Railway, Kolkata, at a cost of , and was sent to Indian Railways. Work is also ongoing to upgrade its old stainless steel third rail to more modern and sustainable aluminium third rail. This could help reduce energy loss by about 84% and solve the problem of voltage drops. This also decrease the time interval between trains to just 90 seconds from 5 minutes. Indian Railways approved the proposal, installation work of Communication Based Train Control signal is expected to be started after conversation of third rail (which is expected to be completed within 2 years) and will be completed within 2–3 years.
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Green Line or Line 2, is the metro corridor to connect Kolkata with Howrah by an underwater tunnel below the Hooghly River. The length was supposed to be , underground and elevated. However, the project was stalled several times due to land acquisition and slum relocation issues. A major route realignment in 2013 increased the length to . The elevated stretch is long while the underground stretch is . The planned intersection with the Blue Line at Central was re-aligned to Esplanade (interchanges with Blue line and Purple line). In September 2019, during the construction of the eastbound tunnel between Sealdah metro station and Esplanade metro station, a Tunnel boring machine hit an Aquifer at Bowbazar, causing a major collapse in the area, delaying work on that section for several months. These issues have caused massive delays to the project, and foreign currency losses had led to an 80 percent cost escalation of the project to nearly .
Between Mahakaran and Howrah, the metro runs under the Hooghly River– the biggest and the only underwater metro tunnel in India. Transfer stations with railways are located at Two major railway stations, Sealdah and Howrah. A new elevated extension from Sector-V to Teghoria was sanctioned a distance of at a budget of in 2016. From Teghoria, passengers can take the Orange Line metro.
The line from Sector-V to Salt Lake Stadium was inaugurated on 13 February 2020 by the then Minister of Railways Piyush Goyal after 11 years of construction. Services to Phoolbagan metro station, the first underground station of the line, were extended on 4 October 2020. The extension added to the existing line. On July 11, 2022, this line was extended till Sealdah. On 6 March 2024, the Esplanade - Howrah Maidan section was inaugurated by Prime Minister Narendra Modi, leaving only 2.9 km between the two functional stations to be joined.
As for a major latest update on 13 January 2025, the remaining 2.9 km between the two functional stations, Esplanade and Sealdah has been successfully connected. As per the details, this initiative will bring down the travel time between the two railway stations to just 11 minutes.
Previously, the stretch from Thakurpukur to Majerhat was surveyed as a branch line of the circular railway, and a metro line from Majerhat to Dakshineswar via Sealdah (interchange with Green line) was planned. This plan was scrapped and a new metro line from further south in Joka to BBD Bagh was sanctioned in 2010–2011 with a total length of at an anticipated cost of . Later the route was truncated to Esplanade. The corridor runs along Diamond Harbour Road, Khidirpur Road, and Jawaharlal Nehru Road, major arterial roads of Kolkata, and has passenger interchange facilities with the Blue Line and Green line at Esplanade and Blue line at Park street. The proposed Esplanade station will not be the same as that of the Blue Line but a different station that will also serve the Green Line. The line now has a new depot in Joka. Due to land acquisition problems and objections from the Ministry of Defence, construction has been delayed several times since the beginning. Defence Ministry objected that the elevated corridor would overlook the Eastern Command headquarters at Fort William, Ordinance Depot at Mominpore. The change in alignment from elevated to underground increased the construction cost of the stretch from to . The work resumed in several phases and new bids were invited by Rail Vikas Nigam Limited (RVNL) in April 2020. It is India's first metro line to run on indigenous head hardened rails, manufactured by Jindal Steel & Power. The extension of this line to IIM and Diamond Park for was sanctioned in the 2012–2013 Budget at a cost of . The work is being executed by RVNL.
The line has 3 phases:
Joka to Majerhat (Phase 1)
Majerhat to Esplanade (Phase 2)
Joka to Diamond Park (Phase 3)
The Mominpur metro station was planned to be built across a 2500 sq. m area. However, the Ministry of Defence objected to the elevated structure, saying that it would overlook the Ordnance depot. This forced RVNL to stall the entire project, and RVNL almost dropped the station from the plan even though it alone would have a projected 20,000 passengers during peak hours. Underground Mominpur station was also not possible due to the sharp gradient from Taratala metro station. After a series of discussions and consultations with the Ministry of Defence and Government of West Bengal in 2016, it was decided to shift the station around northward, near the Alipore Bodyguard Lines. But, after a year Defence Ministry approved the Mominpur metro station in its original location as the change in alignment would have delayed the project and budget overrun. It will be the last elevated station of the corridor. Now, the proposed underground Khidirpur metro station is planned at the Alipore Bodyguard Lines. There were also hurdles regarding clearance for tunneling under defence lands. In 2020, the Defence Ministry eased out the process as lease rent wasn't required anymore for tunneling as long the overground ownership of the land did not change. On 30 December 2022 the Joka - Taratala section and on 6 March 2024, the Taratala - Majerhat section was inaugurated remotely by Honourable Prime Minister Narendra Modi, completing the 7.75 km Phase 1 stretch. On 12 January 2025, it was reported that Metro authorities decided to extend the line to Eden Gardens instead of Esplanade to facilitate commuters and connect Strand Road, BBD Bag and Calcutta High Court. A 1.6 km extension has been planned.
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The work of integrating the Circular Railway from Dum Dum Cantonment to Biman Bandar into a new metro line from Noapara to Netaji Subhash Chandra Bose International Airport was sanctioned in the 2010–2011 budget. The cost of the project is . An eastward extension from Biman Bandar to Barasat over was also sanctioned and included in the 2010–2011 budget. The cost of the project is . The work on this project from Noapara to Barasat is being executed by Metro Railway, Kolkata. Due to multiple delays and hurdles, the total cost of the project had grown to .
Following an objection from the Airports Authority of India (AAI), the route was further reworked. Instead of using the Circular Line's Jessore Road and Biman Bandar railway station, Jessore Road and Jai Hind metro station were planned at-grade and underground, respectively. This stretch will continue underground till Barasat after Prime Minister Narendra Modi's approval which was till New Barrackpore earlier. As of 2024, the construction work has started from the airport to New Barrackpore underground link, and the bidding for the New Barrackpore to Barasat line extension is expected to start.
The Pink Line is the northward extension from Baranagar to Barrackpore []. It was sanctioned at a cost of in the 2010–2011 budget. This line was meant to enable a quick commute from the northernmost suburbs to South Kolkata. The work corridor is being executed by RVNL. As of May 2021, no physical construction has commenced, and the project has been stalled as metro construction would affect the water pipelines along Barrackpore Trunk Road. To avoid this, another proposal was made to continue this line through the Kalyani Expressway. Eleven metro stations were planned on this route.
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A connection between New Garia and Netaji Subhas Chandra Bose International Airport () via EM Bypass, Salt Lake and Rajarhat-New Town was sanctioned to reduce travel time between the southern fringes of Kolkata and the airport. Work on this line was inaugurated by the then Railway Minister Mamata Banerjee on 7 February 2011 with a project deadline of six years. The link between Kavi Subhash and Jai Hind, to be set up at a cost of , will have 24 stations with the terminal Jai Hind metro station being an underground one. The work is executed by Rail Vikas Nigam Limited. Jai Hind metro station will also have a stabling yard, and will be the largest underground facility in the country. This line will have interchanges at Kavi Subhash (with Blue line); Salt Lake Sector V (with Green line) and Teghoria/VIP Road (again with Green line). In July 2020, bids were invited by RVNL to complete the sections left due to various reasons and hurdles.
Initially, the Jai Hind metro station was planned to be elevated. However, the AAI objected that the elevated stretch up to the airport might pose a threat to aircraft, so the route was further reworked and the station was shifted underground, 150 m from the Airport terminal building. As per another revised plan, this line will continue till Barasat and the Yellow line would terminate at Jai Hind. There are also possibilities that Jai Hind metro station would serve as a junction of three lines, i.e. Noapara–Jai Hind, Kavi Subhash–Jai Hind, and Jai Hind–Barasat.
On 6 March 2024, the Kavi Subhash-Hemanta Mukhopadhyay section was inaugurated, completing the 5.4 km Phase 1 stretch.
Proposed expansions
In 2012, RITES, surveyed 16 new routes for connecting the suburban areas to the city. The key routes were:
Majerhat to Ruby via Kalighat and Ballygunge
Basirhat to Tollygunge
New Garia to Haroa via Bhangar
Joka to Mahanayak Uttam Kumar via Thakurpukur
New Garia to Canning via Baruipur along with EM Bypass
Joka to Diamond Harbour along Diamond Harbour Road (Line 3 extension)
Barasat to Barrackpore via SH-2 (Line 4 extension)
Barrackpore to Kalyani via Kalyani Expressway (Line 4/ Line 5 extension)
Madhyamgram to Barrackpore via Sodepur Road and Kalyani Expressway
Branch line of Line 2 from Karunamoyee to Kolkata station
Howrah Maidan – Shalimar – Santragachi (Line 2 extension)
Santragachi to Dhulagarh (Line 2 extension)
Howrah Maidan to Dankuni via Ichapur Road and Benaras Road
Howrah Maidan to Srirampore and Chandannagar via Dankuni, National Highway 2
Howrah Maidan to Belur
Owners and operators
Since the formation of the Metropolitan Transport Project (MTP) in 1969, Kolkata Metro has always been under the Indian Railways, directly or indirectly. It is the only metro in the country to be controlled by Indian Railways. On 29 December 2010, Metro Railway, Kolkata, became the 17th zone of the Indian Railways, completely owned and funded by the Ministry of Railways. Although Kolkata Metro Rail Corporation was formed with 50-50 shares of the Government of West Bengal and the Government of India, as the implementing agency of the East–West Corridor, later majority shares were transferred to Indian Railways. In July 2019, the operation of Green line was handed over to Metro Railway, Kolkata.
Services
Operations
Originally, There are a total of 358 services every day. But, the services and timings were changed due to the COVID-19 pandemic and as of November 2024, it operates between 06:50 and 22:40 IST. Trains operate at an average speed of and stop for about 10 to 20 seconds at each station, depending on the crowd. All stations have display boards showing the terminating station, current time, scheduled time of arrival and estimated time of arrival of trains in Bangla, Hindi and English. Digital countdown clocks are also present in the stations. The coaches of blue and green line have line route-maps and all line have speakers and displays, which provide details of upcoming stations in the three languages. Navigation information is available on Google Maps. Kolkata Metro has launched its own official mobile app 'Metro Ride Kolkata' for android & iOS smartphone users which provides information regarding station, train timing, fare and has online smart card recharge facility along with mobile QR code ticketing.
Seat reservation
In 2008, the Kolkata Metro Railway experimented with the practice of reserving two entire compartments for women. This system was found to be ineffective and caused inconvenience for a lot of commuters (including women) and the plan was dropped.
Now, certain sections of seats in each compartment are reserved for women, senior citizens and the physically challenged. The four-seat sections at each end of a coach are reserved for senior citizens and the physically challenged, and the two middle seat sections, between the general seat sections on each side, are reserved for women.
Fare
The fare is based on the predetermined distance formulas. Kolkata Metro has the lowest starting fare in the country of . For Blue Line, the fare ranges from to ,for Green Line, its to , for Purple Line, the fare ranges from to and Orange line fare range is from to .
Tickets
After using the magnetic ticketing strip system from 1984 to 2011, Kolkata Metro introduced Radio-Frequency Identification (RFID) tokens by Centre for Railway Information Systems (CRIS) in partnership with Keltron in August 2011. The old magnetic strip reader gates were replaced with new RFID readers. The gates are AFC types of gates. These tokens are touched on the machine to enter the station, while to exit from the destination station, it is required to submit the token into the machine. The current tokens are coin-shaped and made of plastic.
QR Code Ticket
The Kolkata Metro has introduced a QR code ticketing system on several lines, including the Blue Line (Line 1), Green Line (Line 2), and Orange Line (Line 6), with plans to extend this system to more lines in the future. This system allows commuters to purchase and use QR code tickets for seamless travel.
Commuters can also opt for mobile QR code tickets via the Metro Ride Kolkata app, available on both Google Play Store and iOS App Store. Only one person can use a mobile ticket at a time. Users have the option to take a screenshot of the QR code ticket for easy access during travel.
Additionally, the app offers a smart card recharge feature, allowing passengers to top up their metro cards digitally. The balance can then be updated at any Add Value Machine (AVM) located in metro stations. The app also provides real-time metro route and timing information, further enhancing the commuting experience.
Smart Card
After introducing RFID tokens, Kolkata metro introduced a Smart Card service provided by CRIS. Earlier, four different types of smart cards were used: Minimum Multi Ride (MMR), Limited Multi Ride (LMR), General Multi Ride (GMR) and Extended Multi Ride (EMR). They were withdrawn on 7 November 2013 and a single type of Smart Card (General Smart Card) was introduced. Two new types of Tourist Smart Cards were also introduced (Tourist Smart Card – I and Tourist Smart Card – II). There is a compulsory refundable security deposit of . The card is common for both the Blue line and the Green line. Online smart card recharge facility was launched on 1 July 2020. These smart cards are not required to be submitted to the AFC gates at the arriving station and can be carried by the passengers. These cards are required to be recharged if the previously recharged money is already spent.
Tourist Smart Card
Two new types of Tourist Smart Cards were also introduced (Tourist Smart Card – I and Tourist Smart Card – II). This type of smart card is for tourists and has unlimited rides. They cost , valid for a day and , valid for three days. A security deposit of is also charged.
Durga Puja special services
The metro railway runs special night-long services during Durga Puja (Maha Saptami to Maha Navami) to help people travel faster and more conveniently for pandal-hopping. The services start at 13:00 and operate till 04:00 the next day. Pre-puja services are also run.
Security
All stations are equipped with closed-circuit cameras, metal detectors and baggage scanners. The Railway Protection Force provides security on the premises. Smoking is strictly prohibited in the metro premises. All stations in the Green Line have half-height and full-height platform screen doors for elevated and underground stations, respectively.
Other facilities
All stations have televisions that broadcast news and songs. WiFi was introduced at Park Street and Maidan metro station in 2016. Gradually, it was expanded to all the stations. The service is provided by Reliance Jio.
Most stations have services such as ATMs, food outlets, and chemist stalls. To ease crowding for recharging smart cards, two Automatic Card recharge machines were installed at Dum Dum. On account of the Swacchota–i–Seba (in English, Cleanliness is service), a nationwide awareness and mobilization campaign on cleanliness, plastic bottle crushers were placed at multiple stations.
Ridership
Kolkata Metro is the 2nd busiest metro system in India. 2,465 travel by every Metro train in Kolkata against 1,110 in Delhi metro. Kolkata Metro carries around 700,000 people daily. The daily and annual ridership has consistently risen since 1984. Low fares and fast and convenient travel have contributed to the high ridership figures. During the 2019 Durga Puja, there was a record ridership of .
Infrastructure
Rolling stock
NGEF rake (Non AC) (decommissioned)
ICF/BHEL rake (Blue Line)
ICF/MEDHA rake (Blue Line, Purple Line, Orange Line)
CRRC Dalian rake (Blue Line)
BEML rake (Green Line)
Depots and yards
There are 5 operational depots now. The Noapara, Tollygunge and New Garia depots serve the Blue Line, the New Garia depot also serve the Orange Line, while the Central Park depot serves the Green Line and the Joka depot serves the Purple Line A depot at New Town for Orange Line and a yard at Airport are under construction.
Stations and electrification
Kolkata Metro has 50 operational stations, of which 21 are underground, 26 are elevated and 3 are at grade. Currently, Noapara is the largest metro station in the system and it will be the interchange station for the Blue Line and Yellow Line. Howrah metro station is the deepest metro station in India. The standard length of platforms in Kolkata Metro is 170 m. The metro stations of Gitanjali and Netaji have the shortest platforms of 163 m. The average length between any two stations is . The shortest distance is between Central and Chandni Chowk, and the longest distance is between Noapara and Baranagar. Since the Kolkata Metro has electrification, electricity substations were built in Jatin Das Park, Central and Shyambazar.
Signalling and telecommunication
Blue Line: Trains operate on typical Indian Railways automatic signaling technology. A Route Relay Interlocking System has been provided at New Garia depot and Tollygunge depot and Electronic Interlocking has been provided at Noapara depot to facilitate the prompt withdrawal and injection of rakes and to perform shunting operations inside the car shed for maintenance purposes. The Train Protection and Warning System (TPWS) is provided throughout the Metro Railway. It is designed to prevent collisions caused by human (operator) error. A Train Describer System and Auto Train Charting are utilized to assist the operation control center in monitoring and planning train movements in real-time. An Integrated Power Supply System and microprocessor-based Data Logger System have also been provided. An integrated system of STM-1 and STM-4 optical fibre cable is used for all telecommunication, signaling, SCADA and other circuits in Blue line. The service is provided by RailTel.
An upgrade of the existing signaling system of the North–south corridor from Indian Railways Signalling to Communications-based train control (CBTC) was planned by Metro Railway, Kolkata at and the proposal was sent to Indian Railways, so that time interval between trains can be decreased to just 90 seconds from 5 minutes. In August 2019, Indian Railways gave the go-ahead to the proposal, and installation work is supposed to be complete within 2–3 years.
Green Line: Unlike the previous line, the Green line adopted a more advanced CBTC system. It has cab signalling and a centralised automatic train control system consisting of automatic operation, protection, and signaling modules. The signaling system is provided by Italy-based company Ansaldo STS. The other signalling equipment includes an integrated system with fibre optic cable, SCADA, radios, and a public-address system.
Purple and Orange Lines: CBTC systems are currently being installed in both the Purple and Orange Lines.
Public address system
PA systems are present at all stations and their premises. A station master can make a necessary announcement to the passengers and staff, overriding the ongoing local announcement. Train PA systems are controlled by the motormen for announcements to passengers on the particular train.
Issues
Since the Kolkata Metro was constructed in the 1970s, there were some technical limitations. Due to the tunnel dimensions, and being under Indian Railways, Kolkata Metro opted for an Indian metre gauge shell (2.7 m width) mounted upon broad gauge bogies. The rakes have to be custom built and require a special assembly line involving additional costs thus limiting the options for rake manufacturers for Blue line. From its inception, the coaches were manufactured by ICF, which lacked the pre-requisite knowledge for manufactured non-air-conditioning rakes. The 3000 and 4000 series rakes were faulty and delivered without any trials. In addition, Indian Railways signaling is used instead of European signaling. All of these factors have led to snags, delays, and accidents.
Unlike Delhi Metro, Kolkata Metro is owned and operated by Indian Railways instead of an autonomous body, and it relies solely on Indian Railways for every decision, from funding to route realignment.
Underground tunnels
In densely populated areas, there are no free spaces left to build elevated metro tracks and stations. As a result, underground systems are constructed in these areas. However, the construction of an underground metro tunnel in the Bowbazar area has caused cracks in the houses of many residents. Consequently, the metro authority had to evacuate the people, leading to massive delays and slow development in the construction of the metro railway in that area.
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33067333 | https://en.wikipedia.org/wiki/Splinter | Splinter | A splinter (also known as a sliver) is a fragment of a larger object, or a foreign body that penetrates or is purposely injected into a body. The foreign body must be lodged inside tissue to be considered a splinter. Splinters may cause initial pain through ripping of flesh and muscle, or infection through bacteria on the foreign object.
Splinters are primarily made of wood, but there are many other types, for example, other common types of splinters are, hair, glass, plastic, metal, and spines of animals.
As with any wound that breaks the skin, splinters can lead to infection, which if left untreated could develop into more serious complications. If a splinter is in the body for more than 2–3 days, or if the wound shows signs of inflammation or tenderness (whether the splinter was removed or not), advice should be sought from a doctor.
Getting a splinter
Generally, a splinter causes an initial feeling of pain as the sharp object makes its initial penetration through the body. Through this penetration, the object cuts through the cutaneous layer of the skin, and settles in the subcutaneous layer of the skin, and can even penetrate further down, breaking the sub-cutaneous layer, settling in muscle tissue, or even the bone. Some splinters will remain in place, but most will continue to migrate through the body (eg. hair splinters), further damaging their surroundings.
Types
According to the American Academy of Family Physicians, the most common foreign bodies contracted by people fall into two official classes: biological splinters, and nonbiological splinters. In the biological class, splinters include bone, fish spines, teeth, hair and wood. Although, in the nonbiological class, common splinters contracted are glass, metal, aluminum, fishhooks, pencil graphite, and plastic.
Rarely, people may become infected with splinters from more unusual sources. Common cases of exotic foreign bodies include sea urchins, insect stings, stingray spines, and even grenade shrapnel.
Specific details of some splinters
Wood: this type of splinter is contracted from lumber or other vegetative materials. Wood splinters must be removed from wounds because they are associated with inflammation and risk of infection. Larger or deeper splinters can result in difficult removal, or localization of the foreign body.
Fishhooks: fishhooks that become lodged in the skin are problematic because of the barbs found on the ends of most fishhooks. These barbs are designed to make removal difficult, and if care is not taken, the victim can experience tearing of not only the flesh, but the muscle as well. The most common injuries caused by fishhooks occur in the hand, face, scalp, foot, and eye.
Glass: One study found that patients were more likely to feel sensations from glass present in their skin than from any other kind of splinter. Though glass is generally detectable by radiography and is radiopaque, there is limited ability for radiography to detect glass fragments smaller than 2 mm. Most glass splinters are inert, and generally lack the ability to migrate to other regions of the body.
Other: Pencil lead and other graphite foreign bodies, once lodged in the cutaneous layer of the skin, can cause permanent pigment tattooing if not removed immediately. Metallic bodies range from BB pellets to grenade shrapnel. Smaller objects can be removed without much difficulty if the depth of the wound remains superficial, but if the wound does not protrude past the subcutaneous layers of the skin, and remains inert, the object can actually remain in place. In larger objects, fragments that remain superficial in one's body may be removed without much trouble, but if wounds protrude past the subcutaneous layers of the skin and even into the muscular area or near vital organs, such objects must be left alone and immediate medical attention must be sought.
Hair: Short lengths of hair, especially stiff hair such as trimmed beard hair or pet hair, can work their way under the skin of the feet or hands. Hair splinters (cutaneous pili migrans) are commonly experienced by hairdressers and dog groomers. This is distinct from an ingrown hair, where a hair still attached to its follicle grows back under (or fails to emerge from) the skin. Furthermore, unlike an ingrown hair, the hair that is embedded by a hair splinter may not be one's own; the hair may have either belonged to another person or to an animal.
Detection
Splinters are often first detected by the person with the splinter in their body. There are many signs that a splinter has entered one's body.
Signs of a hidden foreign body
Puncture wound
Blood-stained injury track of a fresh wound
Sharp pain with deep palpation over a puncture wound
Discoloration beneath the epidermis
Wound that elicits pain with movement
Wound that fails to heal
Abscess (with sterile culture)
Pain associated with a mass
Mass under the epidermis
Chronically draining purulent wound
Cyst
Granuloma formation
Sterile monoarticular arthritis
Periosteal reactions
Osteomyelitis
Pseudotumors of bone
Delayed tendon or nerve injury
Imaging
If manual detection and localization fail, the main methods for medical imaging of splinters are:
Projectional radiography – used to locate bone, fish spines, glass, gravel-stone, metal, aluminum, pencil graphite, some plastics, teeth, and some wood (e.g., spines, cactus, thorns)
Medical ultrasonography – used to locate glass, metal, pencil graphite, some plastics, stone, and some types of wood.
Small wooden splinters (1–4 mm) distant from bones are most easily detected by ultrasonography, while CT scan and magnetic resonance imaging have higher sensitivity for those near bones.
Removal
There are many medical techniques to remove splinters safely. Common medical techniques include the Elliptical Technique and the String Technique. In the elliptical technique the surrounding area of the splinter is sliced in an elliptical formation. From there the flesh in the elliptical area is cut (in the shape of an upside-down cone) and the whole chunk of flesh containing the splinter is removed. The Needle Cover Technique is limited to fishhook removal. A string is looped around the base of the hook, and as the hook is pressed further into the skin, the string is pulled, allowing the barbs to be unhooked from muscle and follow the path of the rest of the hook out of the body without snagging any additional flesh.
Since the splinter has penetrated through a physical barrier of the body it allows for an individual to get an infection. The opening from the splinter will make it easier for bacteria to get into the body. It is strongly encouraged for the removal of a splinter before falling victim to an infection.
Infection
Infection is usually determined by the duration of time that the foreign object remains lodged in the human body. Objects that have included poison, deep penetration, dirt, or bite injuries generally result in a shorter time until infection is notable. According to the AAFP, patients that are older, or have diabetes, or have wounds that are longer, wider, more jagged or deeper, have a much higher risk of infection. Simply the easiest way to avoid infection is to completely remove the splinters or foreign body as soon as possible.
Though infection is generally the largest complication encountered with splinters, ranging from 1.1 to 12 percent presence, the use of antibiotics in non-bite cases is generally deemed unnecessary by the medical community. Though cases are rare, infection of foreign body wounds can result in cases of tetanus.
One case of tetanus contraction through a splinter was seen in Ohio in 1993. An 80-year-old woman was presented to an ED with dysphagia and a stiff jaw. Not long after a preliminary checkup, a wood splinter was found to have been lodged in her chin for approximately 1 week; the area was erythematous with active purulent drainage. The woman was diagnosed with tetanus, admitted to the hospital, and begun on a regimen of 3,000 units of tetanus immune globulin, tetanus toxoid, and intravenous clindamycin. Despite aggressive treatment, including assisted mechanical ventilation, the patient died 15 days later from the effects of her primary infection. The woman had no history of previous tetanus vaccinations despite previous care for a wound and ongoing medical attention for hypertension.
Since most splinters are made of organic matter, they are much more dangerous than other types of things puncturing the body. Splinters are usually infected with many bacteria which then turn into an infection such as tetanus. Due to a splinter being made of organic matter, it makes it much more difficult for the body to get rid of it.
| Biology and health sciences | Types | Health |
1423441 | https://en.wikipedia.org/wiki/Free%20range | Free range | Free range denotes a method of farming husbandry where the animals can roam freely outdoors for at least part of the day, rather than being confined in an enclosure for 24 hours each day.
On many farms, the outdoors ranging area is fenced, thereby technically making this an enclosure, however, free range systems usually offer the opportunity for the extensive locomotion and sunlight that is otherwise prevented by indoor housing systems. Free range may apply to meat, eggs or dairy farming.
The term is used in two senses that do not overlap completely: as a farmer-centric description of husbandry methods, and as a consumer-centric description of them. There is a diet where the practitioner only eats meat from free-range sources called ethical omnivorism.
In ranching, free-range livestock are permitted to roam without being fenced in, as opposed to intensive animal farming practices such as the concentrated animal feeding operation. In many agriculture-based economies, free-range livestock are quite common.
History
If one allows "free range" to include "herding", free range was a typical husbandry method at least until the development of barbed wire and chicken wire. The generally poor understanding of nutrition and diseases before the twentieth century made it difficult to raise many livestock species without giving them access to a varied diet, and the labor of keeping livestock in confinement and carrying all their feed to them was prohibitive except for high-profit animals such as dairy cattle.
In the case of poultry, free range was the dominant system until the discovery of vitamins A and D in the 1920s, which allowed confinement to be practised successfully on a commercial scale. Before that, green feed and sunshine (for the vitamin D) were necessary to provide the necessary vitamin content. Some large commercial breeding flocks were reared on pasture into the 1950s. Nutritional science resulted in the increased use of confinement for other livestock species in much the same way.
United States
In the United States, the USDA free range regulations currently apply only to poultry and indicate that the animal has been allowed access to the outside. To be considered "free range," the poultry must have access to the outdoors for more than 51% of the animal's life. However, USDA regulations make no mention of the quality or size of the outside range. The Certified Humane Program offers third-party certification for producers who meet minimum standards, including providing access to grass pastures, traditional nests, and "dust areas to perform natural behaviors".
The term "free range" is mainly used as a marketing term rather than a husbandry term, meaning something on the order of, "low stocking density", "pasture-raised", "grass-fed", "old-fashioned", "humanely raised", etc.
There have been proposals to regulate USDA labeling of products as free range within the United States. what constitutes raising an animal "free range" is almost entirely decided by the producer of that product, and is frequently inconsistent with consumer ideas of what the term means.
Free-range poultry
In poultry-keeping, "free range" is widely confused with yarding, which means keeping poultry in fenced yards. Yarding, as well as floorless portable chicken pens ("chicken tractors") may have some of the benefits of free-range livestock but, in reality, the methods have little in common with the free-range method.
A behavioral definition of free range is perhaps the most useful: "chickens kept with a fence that restricts their movements very little." This has practical implications. For example, according to Jull, "The most effective measure of preventing cannibalism seems to be to give the birds good grass range." De-beaking was invented to prevent cannibalism for birds not on free range, and the need for de-beaking can be seen as a litmus test for whether the chickens' environment is sufficiently "free-range-like".
The U.S. Department of Agriculture Food Safety and Inspection Service (FSIS) requires that chickens raised for their meat have access to the outside in order to receive the free-range certification. However, what qualifies as "access" is not defined. There is no requirement for access to pasture, and there may be access to only dirt or gravel. Free-range chicken eggs, however, have no legal definition in the United States. Likewise, free-range egg producers have no common standard on what the term means.
The broadness of "free range" in the U.S. has caused some people to look for alternative terms. "Pastured poultry" is a term promoted by farmer/author Joel Salatin for broiler chickens raised on grass pasture for all of their lives except for the initial brooding period. The Pastured Poultry concept is promoted by the American Pastured Poultry Producers' Association (APPPA), an organization of farmers raising their poultry using Salatin's principles. This term is not defined by the USDA and has no legal definition. To use a term like pasture-raised, producers must submit documentation to the FSIS of continuous free access to the outdoors, but do not need to include additional terminology to define pasture-raised on its label.
Free-range livestock
Traditional American usage equates "free range" with "unfenced", and with the implication that there was no herdsman keeping them together or managing them in any way. Legally, a free-range jurisdiction allowed livestock (perhaps only of a few named species) to run free, and the owner was not liable for any damage they caused. In such jurisdictions, people who wished to avoid damage by livestock had to fence them out; in others, the owners had to fence them in.
The USDA has no specific definition for "free-range" beef, pork, and other non-poultry products. All USDA definitions of "free-range" refer specifically to poultry.
In a December 2002 Federal Register notice and request for comments (67 Fed. Reg. 79552), USDA's Agricultural Marketing Service proposed "minimum requirements for livestock and meat industry production/marketing claims". Many industry claim categories are included in the notice, including breed claims, antibiotic claims, and grain fed claims. "Free Range, Free Roaming, or Pasture Raised" would be defined as "livestock that have had continuous and unconfined access to pasture throughout their life cycle" with an exception for swine ("continuous access to pasture for at least 80% of their production cycle"). In a May 2006 Federal Register notice (71 Fed. Reg. 27662), the agency presented a summary and its responses to comments received in the 2002 notice, but only for the category "grass (forage) fed" which the agency stated was to be a category separate from "free range". Comments received for other categories, including "free range", are to be published in future Federal Register editions.
European Union
The European Union regulates marketing standards for egg farming which specifies the following (cumulative) minimum conditions for the free-range method:
hens have continuous daytime access to open-air runs, except in the case of temporary restrictions imposed by veterinary authorities,
the open-air runs to which hens have access is mainly covered with vegetation and not used for other purposes except for orchards, woodland and livestock grazing if the latter is authorised by the competent authorities,
the open-air runs must at least satisfy the conditions specified in Article 4(1)(3)(b)(ii) of Directive 1999/74/EC whereby the maximum stocking density is not greater than 2500 hens per hectare of ground available to the hens or one hen per at all times and the runs are not extending beyond a radius of from the nearest pophole of the building; an extension of up to from the nearest pophole of the building is permissible provided that a sufficient number of shelters and drinking troughs within the meaning of that provision are evenly distributed throughout the whole open-air run with at least four shelters per hectare.
Otherwise, egg farming in EU is classified into 4 categories: Organic (ecological), Free Range, Barn, and Cages.) The mandatory labelling on the egg shells attributes a number (which is the first digit on the label) to each of these categories: 0 for Organic, 1 for Free Range, 2 for Barn and 3 for Cages.
There are EU regulations about what free-range means for laying hens and broilers (meat chickens) as indicated above. However, there are no EU regulations for free-range pork, so pigs could be indoors for some of their lives. In order to be classified as free-range, animals must have access to the outdoors for at least part of their lives.
United Kingdom
Pigs:
Free-range pregnant sows are kept in groups and they are often provided with straw for bedding, rooting and chewing. Around 40% of UK sows are kept free-range outdoors and farrow in huts on their range.
Egg laying hens:
Cage-free egg production includes barn, free-range and organic systems. The UK is the largest free-range egg producer in the Europe. Free-range systems are the most popular of the non-cage alternatives, accounting for around 57% of all eggs, compared to 2% in barns and 2% organic. In free-range systems, hens are housed to a similar standard as the barn or aviary.
Free-range rearing of pullets:
Free range rearing of pullets for egg-laying is now being pioneered in the UK by various poultry rearing farms. In these systems, the pullets are allowed outside from as young as 4 weeks of age, rather than the conventional systems where the pullets are reared in barns and allowed out at 16 weeks of age
Meat chickens:
Free-range broilers are reared for meat and are allowed access to an outdoor range for at least 8 hours each day. RSPCA standards state that in order for chickens to be free range, there must not be more than 13 chickens per square meter. Free-range broiler systems use slower-growing breeds of chicken to improve welfare, meaning they reach slaughter weight at 16 weeks of age rather than 5–6 weeks of age in standard rearing systems.
Turkeys:
Free-range turkeys have continuous access to an outdoor range during the daytime. The range should be largely covered in vegetation and allow more space. Access to fresh air and daylight means better eye and respiratory health. The turkeys are able to exercise and exhibit natural behaviour resulting in stronger, healthier legs. Free-range systems often use slower-growing breeds of turkey.
Free range dairy:
Farms supplying milk under the free range dairy brand abide by the pasture promise, meaning the cows will have access to pasture land to graze for a minimum of 180 days and nights a year. There is evidence to suggest that milk from grass contains higher levels of fats such as omega-3 and conjugated linoleic acid (CLA). Additionally free range dairy is giving consumers more choice as to where their milk comes from. Free range dairy provides the consumer with reassurance that the milk they drink has come from cows with the freedom to roam and can graze in their natural habitat.
Australia
Australian standards in relation to free-range production are largely espoused in third-party certification trade marks due to the absence of any significant legally binding legislation. A number of certification bodies are utilised by rearers to identify their products with a particular level of animal welfare standards. In events where producers do not choose to use a certified trade mark and merely state that their product is 'free range', the producer is bound by consumer expectations and perceptions of what constitutes free range. Producers are generally thought to be bound to Model Codes of Practice of Animal Welfare published by the CSIRO, and in some states this forms part of legislation.
Egg laying hens
In Australia, three farming methods for the production of eggs are utilised. In 2011, traditional cage (or battery) eggs accounted for 42% of value, barn-laid eggs account for 10% of value, free-range eggs accounted for 44% of value, and organic eggs accounted for 4% of value. Increased demand for free range eggs due to customer concerns over animal welfare has led to a number of different standards developing in relation to three core welfare measures – indoor stocking density, outdoor stocking density, and beak trimming. The Model Code of Practice recommends practices for free range farming with the following standards:
Maximum stocking densities indoors of 30 kg/m2, equivalent to about 14–15 birds per square metre.
Maximum outdoor stocking density of 1500 birds per hectare, although this can be increased with rotation onto fresh pasture
Access to the outdoor range for a minimum of 8 hours per day, except in adverse weather conditions
2 metres worth of popholes per 1,000 birds for access to the range
Beak trimming is permitted, and to be undertaken by an accredited operator
The above standards are not always met, and on some occasions producers may want more ethical standards. As such, certified trade marks play a significant role in the determining of what constitutes free range. The key certifications used for layer hens in Australia include the following:
Egg Corp Assured is the weakest standard, set by the industry peak group and largely based on the Model Code of Practice. Egg Corp Assured differs in that it interprets the outdoor stocking density figure as largely irrelevant to welfare. Egg Corp Assured has been known to certify farms running up to 44,000 birds per hectare outdoors, far in excess of recommendations. Like the Model Code of practice, beak trimming is allowed and indoor densities run up to 15 birds per square metre.
RSPCA Approved Farming is a standard that can be applied to both barn-laid and free-range egg producers. Farms using this certification must have an indoor density of 9 birds per square metre indoors on slats, or 7 birds per square metre indoors in a deep-litter system. The standards dictate a maximum outdoor density of 1500 per hectare without rotation, or 2500 birds per hectare with rotation, and beak trimming is allowed.
Free Range Egg & Poultry Australia (FREPA) standards provides a sliding scale for indoor density, with 10 birds per square metre allowed only in enclosures housing less than 1000 birds, and 6 birds per square metre the maximum for barns with over 4000 birds. Nothing is said in the standards about outdoor density, thus it is assumed that farmers must meet the standards under the Model Code. Beak trimming is allowed under this certification.
Humane Choice True Free Range standards are some of the most sound as far as animal welfare is concerned. Beak trimming or any other mutilation is not permitted, perches must be provided, and maximum flock numbers cannot be greater than 2500 per barn. The outdoor stocking density is 1500 birds per hectare, and the indoor density is 5 birds per square metre.
Australian Certified Organic Standards include criteria on feed content and the use of pesticides in addition to animal welfare requirements. The indoor density is a maximum of 8 birds per square metre, although most operators under this standard list their density as 5 birds per square metre. The outdoor density is 1000 birds per hectare, and beak trimming is not permitted.
Chicken meat
In Australia, free range and organic chicken accounts for about 16.6% of value in the poultry market. This percentage is expected to grow to up to 25% in the next 5 years. No meat birds are raised in cages in Australia. There are three main certification trademarks in this market.
Free Range Egg & Poultry Australia (FREPA) standards are those in which most supermarket brands of free range chicken meat are accredited under. These standards require indoor stocking densities of up to 30 kg/m2 indoors (about 15 birds per square metre), and beak trimming is not permitted. Outdoor stocking density is not stated, but it is understood that the outdoor range must be at a minimum 1.5 times the floor area of inside the barn.
RSPCA Approved Farming standards for free range require an indoor stocking density of about 17 birds per square metre, and outdoor densities of up to 17 birds per square metre. No beak trimming is allowed under this system.
Australian Certified Organic standards dictate a maximum indoor stocking density of up to 12 birds per square metre indoors, and 2500 birds per hectare outdoors. These standards require perches, and prevent large, conventional broiler sheds.
| Technology | Agriculture_2 | null |
1424029 | https://en.wikipedia.org/wiki/Fucus | Fucus | Fucus is a genus of brown algae found in the intertidal zones of rocky seashores almost throughout the world.
Description and life cycle
The thallus is perennial with an irregular or disc-shaped holdfast or with haptera. The erect portion of the thallus is dichotomous or subpinnately branched, flattened and with a distinct midrib. Gas-filled pneumatocysts (air-vesicles) are present in pairs in some species, one on either side of the midrib. The erect portion of the thallus bears cryptostomata and caecostomata (sterile surface cavities). The base of the thallus is stipe-like due to abrasion of the tissue lateral to the midrib and it is attached to the rock by a holdfast. The gametangia develop in conceptacles embedded in receptacles in the apices of the final branches. They may be monoecious or dioecious.
These algae have a relatively simple life cycle and produce only one type of thallus which grows to a maximum size of 2 m. Fertile cavities, the conceptacles, containing the reproductive cells are immersed in the receptacles near the ends of the branches. After meiosis, oogonia and antheridia, the female and male reproductive organs, produce egg cells and sperm respectively that are released into the sea where fertilisation takes place. The resulting zygote develops directly into the diploid plant. This contrasts with the life cycle of the flowering plant, where the egg cells and sperm are produced by a haploid multicellular generation, albeit very strongly reduced, and the egg cells are fertilised within the ovules of the parent plant and then released as seeds.
Distribution and ecology
Species of Fucus are recorded almost worldwide. They are dominant on the shores of the British Isles, the northeastern coast of North America and California.
In the British Isles these larger brown algae occur on sheltered shores in fairly well defined zones along the shore from high-water mark to below low water mark. On the more exposed shores not all of these species can be found and on very exposed shores few, if any, occur. Pelvetia canaliculata forms a zone at the top of the shore. Just below this Fucus spiralis, Fucus vesiculosus, Fucus serratus and Laminaria form clear zones, one below the other, along the shore down to low water mark. On sheltered shores Ascophyllum nodosum usually forms a broad and dominating zone along the shore at the mid-littoral. Other brown algae can be found at the low-littoral such as Himanthalia, Laminaria saxatilis and Alaria esculenta. Small green and red algae and animals occur, protected under these large brown algae. When washed up on beaches, kelp flies such as Coelopa pilipes feed and breed on Fucus algae.
Uses
In Scotland and Norway, up until the mid-19th century, several seaweed species from Fucus and other genera were harvested, dried, burned to ash, and further processed to become "kelp", which was a type of soda ash that was less costly in Britain than the barilla imported from Spain. It has an alkali content of about 2.5%–5% that was mainly sodium carbonate (Na2CO3), used in soapmaking, glassmaking, and other industries. The purest barilla had a sodium carbonate concentration of about 30%. The seaweed was also used as fertilizer for crop land in the same areas in which it was harvested.
Fucus species can also be used for thalassotherapy, along with other species such as Turkish towel (Chondracanthus exasperatus), feather boa (Egregia menziesii), and finger kelp (Laminaria digitata).
In 2005, it was announced that bacteria grown on Fucus have the ability to attack and kill the MRSA superbacterium.
Because of their easily accessible apolar eggs and free-floating zygotes, several Fucus species have been used as model organisms to study cell polarity, the development of growth axes, and the role of the cell wall in establishing and maintaining cell identity.
Predator
The seaweed fly, Coelopa frigida, together with other species of Coelopa, are known to feed, mate, and create habitats out of different species of Fucus. This is of particular notice when the Fucus is stranded on the beach and not when it is submerged under seawater. With increasing amounts of seaweed washing up on shores, there is an increasing recognition of Fucus and their close pairing with Coelopa.
Taxonomy
This list of species of Fucus excludes names of uncertain status:
Fucus atomarius (Woodward) Bertoloni
Fucus ceranoides L. – horned wrack *
Fucus chalonii Feldmann
Fucus cottonii M. J. Wynne & Magne *(=Fucus cottonii M.J.Wynne & Magne nom. illeg.)
Fucus distichus L. *
Fucus evanescens C.Agardh *
Fucus furcatus Stackhouse, 1801
Fucus guiryi G. I. Zardi, K. R. Nicastro, E. A. Serrão, G. A. Pearson *
Fucus gardneri P. C. Silva
Fucus lagasca Clemente, 1807
Fucus mytili Nienburg
Fucus nereideus Lightfoot
Fucus radicans L. Bergström & L. Kautsky, 2005
Fucus serratus L. – toothed wrack *
Fucus spermophorus L.
Fucus spiralis L. – spiral wrack *
Fucus tendo L.
Fucus vesiculosus L. – bladder wrack *
Fucus virsoides J. Agardh
* Species recorded around the coast of Britain.
Fucus distichus
F. distichus is up to 10 cm long with a short stout cylindrical stipe, branching dichotomous, flat and with a mid-rib. F. distichus subsp. edentatus was first described from Shetland by Börgesen in 1903. Powell found F. distichus subsp. anceps on the north coast of Caithness. It had also been recorded from: Orkney, Fair Isle, St Kilda and the Outer Hebrides in Scotland; in Ireland it had been recorded from Counties Clare, Donegal and Kerry. Two subspecies of F. distichus (subsp. anceps and subsp. edentatus) have been described from the British Isles.
Fucus distichus is the organism used as a model to study the development of cell polarity, since it forms an apolar zygote that can develop polarity given a varying number of gradients.
Fucus serratus
F. serratus, toothed wrack, is the most distinctive of all the Fucus species. It clearly shows a distinctive serrated edge to the frond not shown by the other species of the genus.
Fucus spiralis
F. spiralis is one of the three most common algae on the shores of the British Isles. It grows to about 40 cm long and does not show air bladders as found on F.vesiculosus or toothed edges as found on F. serratus. It forms a zone near the top of the shore above the zones of F. vesiculosus and F. serratus.
Fucus vesiculosus
This is one of the most common species of Fucus, common on most shores in the mid-littoral. It has the common name "bladder wrack", and is readily identified by a distinct mid-rib and air vesicles in pairs on either side of the mid-rib.
| Biology and health sciences | SAR supergroup | Plants |
1424309 | https://en.wikipedia.org/wiki/Differential%20equation | Differential equation | In mathematics, a differential equation is an equation that relates one or more unknown functions and their derivatives. In applications, the functions generally represent physical quantities, the derivatives represent their rates of change, and the differential equation defines a relationship between the two. Such relations are common in mathematical models and scientific laws; therefore, differential equations play a prominent role in many disciplines including engineering, physics, economics, and biology.
The study of differential equations consists mainly of the study of their solutions (the set of functions that satisfy each equation), and of the properties of their solutions. Only the simplest differential equations are solvable by explicit formulas; however, many properties of solutions of a given differential equation may be determined without computing them exactly.
Often when a closed-form expression for the solutions is not available, solutions may be approximated numerically using computers, and many numerical methods have been developed to determine solutions with a given degree of accuracy. The theory of dynamical systems analyzes the qualitative aspects of solutions, such as their average behavior over a long time interval.
History
Differential equations came into existence with the invention of calculus by Isaac Newton and Gottfried Leibniz. In Chapter 2 of his 1671 work Methodus fluxionum et Serierum Infinitarum, Newton listed three kinds of differential equations:
In all these cases, is an unknown function of (or of and ), and is a given function.
He solves these examples and others using infinite series and discusses the non-uniqueness of solutions.
Jacob Bernoulli proposed the Bernoulli differential equation in 1695. This is an ordinary differential equation of the form
for which the following year Leibniz obtained solutions by simplifying it.
Historically, the problem of a vibrating string such as that of a musical instrument was studied by Jean le Rond d'Alembert, Leonhard Euler, Daniel Bernoulli, and Joseph-Louis Lagrange. In 1746, d’Alembert discovered the one-dimensional wave equation, and within ten years Euler discovered the three-dimensional wave equation.
The Euler–Lagrange equation was developed in the 1750s by Euler and Lagrange in connection with their studies of the tautochrone problem. This is the problem of determining a curve on which a weighted particle will fall to a fixed point in a fixed amount of time, independent of the starting point. Lagrange solved this problem in 1755 and sent the solution to Euler. Both further developed Lagrange's method and applied it to mechanics, which led to the formulation of Lagrangian mechanics.
In 1822, Fourier published his work on heat flow in Théorie analytique de la chaleur (The Analytic Theory of Heat), in which he based his reasoning on Newton's law of cooling, namely, that the flow of heat between two adjacent molecules is proportional to the extremely small difference of their temperatures. Contained in this book was Fourier's proposal of his heat equation for conductive diffusion of heat. This partial differential equation is now a common part of mathematical physics curriculum.
Example
In classical mechanics, the motion of a body is described by its position and velocity as the time value varies. Newton's laws allow these variables to be expressed dynamically (given the position, velocity, acceleration and various forces acting on the body) as a differential equation for the unknown position of the body as a function of time.
In some cases, this differential equation (called an equation of motion) may be solved explicitly.
An example of modeling a real-world problem using differential equations is the determination of the velocity of a ball falling through the air, considering only gravity and air resistance. The ball's acceleration towards the ground is the acceleration due to gravity minus the deceleration due to air resistance. Gravity is considered constant, and air resistance may be modeled as proportional to the ball's velocity. This means that the ball's acceleration, which is a derivative of its velocity, depends on the velocity (and the velocity depends on time). Finding the velocity as a function of time involves solving a differential equation and verifying its validity.
Types
Differential equations can be divided into several types. Apart from describing the properties of the equation itself, these classes of differential equations can help inform the choice of approach to a solution. Commonly used distinctions include whether the equation is ordinary or partial, linear or non-linear, and homogeneous or heterogeneous. This list is far from exhaustive; there are many other properties and subclasses of differential equations which can be very useful in specific contexts.
Ordinary differential equations
An ordinary differential equation (ODE) is an equation containing an unknown function of one real or complex variable , its derivatives, and some given functions of . The unknown function is generally represented by a variable (often denoted ), which, therefore, depends on . Thus is often called the independent variable of the equation. The term "ordinary" is used in contrast with the term partial differential equation, which may be with respect to more than one independent variable.
Linear differential equations are the differential equations that are linear in the unknown function and its derivatives. Their theory is well developed, and in many cases one may express their solutions in terms of integrals.
Most ODEs that are encountered in physics are linear. Therefore, most special functions may be defined as solutions of linear differential equations (see Holonomic function).
As, in general, the solutions of a differential equation cannot be expressed by a closed-form expression, numerical methods are commonly used for solving differential equations on a computer.
Partial differential equations
A partial differential equation (PDE) is a differential equation that contains unknown multivariable functions and their partial derivatives. (This is in contrast to ordinary differential equations, which deal with functions of a single variable and their derivatives.) PDEs are used to formulate problems involving functions of several variables, and are either solved in closed form, or used to create a relevant computer model.
PDEs can be used to describe a wide variety of phenomena in nature such as sound, heat, electrostatics, electrodynamics, fluid flow, elasticity, or quantum mechanics. These seemingly distinct physical phenomena can be formalized similarly in terms of PDEs. Just as ordinary differential equations often model one-dimensional dynamical systems, partial differential equations often model multidimensional systems. Stochastic partial differential equations generalize partial differential equations for modeling randomness.
Non-linear differential equations
A non-linear differential equation is a differential equation that is not a linear equation in the unknown function and its derivatives (the linearity or non-linearity in the arguments of the function are not considered here). There are very few methods of solving nonlinear differential equations exactly; those that are known typically depend on the equation having particular symmetries. Nonlinear differential equations can exhibit very complicated behaviour over extended time intervals, characteristic of chaos. Even the fundamental questions of existence, uniqueness, and extendability of solutions for nonlinear differential equations, and well-posedness of initial and boundary value problems for nonlinear PDEs are hard problems and their resolution in special cases is considered to be a significant advance in the mathematical theory (cf. Navier–Stokes existence and smoothness). However, if the differential equation is a correctly formulated representation of a meaningful physical process, then one expects it to have a solution.
Linear differential equations frequently appear as approximations to nonlinear equations. These approximations are only valid under restricted conditions. For example, the harmonic oscillator equation is an approximation to the nonlinear pendulum equation that is valid for small amplitude oscillations.
Equation order and degree
The order of the differential equation is the highest order of derivative of the unknown function that appears in the differential equation.
For example, an equation containing only first-order derivatives is a first-order differential equation, an equation containing the second-order derivative is a second-order differential equation, and so on.
When it is written as a polynomial equation in the unknown function and its derivatives, its degree of the differential equation is, depending on the context, the polynomial degree in the highest derivative of the unknown function, or its total degree in the unknown function and its derivatives. In particular, a linear differential equation has degree one for both meanings, but the non-linear differential equation is of degree one for the first meaning but not for the second one.
Differential equations that describe natural phenomena almost always have only first and second order derivatives in them, but there are some exceptions, such as the thin-film equation, which is a fourth order partial differential equation.
Examples
In the first group of examples u is an unknown function of x, and c and ω are constants that are supposed to be known. Two broad classifications of both ordinary and partial differential equations consist of distinguishing between linear and nonlinear differential equations, and between homogeneous differential equations and heterogeneous ones.
Heterogeneous first-order linear constant coefficient ordinary differential equation:
Homogeneous second-order linear ordinary differential equation:
Homogeneous second-order linear constant coefficient ordinary differential equation describing the harmonic oscillator:
Heterogeneous first-order nonlinear ordinary differential equation:
Second-order nonlinear (due to sine function) ordinary differential equation describing the motion of a pendulum of length L:
In the next group of examples, the unknown function u depends on two variables x and t or x and y.
Homogeneous first-order linear partial differential equation:
Homogeneous second-order linear constant coefficient partial differential equation of elliptic type, the Laplace equation:
Homogeneous third-order non-linear partial differential equation, the KdV equation:
Existence of solutions
Solving differential equations is not like solving algebraic equations. Not only are their solutions often unclear, but whether solutions are unique or exist at all are also notable subjects of interest.
For first order initial value problems, the Peano existence theorem gives one set of circumstances in which a solution exists. Given any point in the xy-plane, define some rectangular region , such that and is in the interior of . If we are given a differential equation and the condition that when , then there is locally a solution to this problem if and are both continuous on . This solution exists on some interval with its center at . The solution may not be unique. (See Ordinary differential equation for other results.)
However, this only helps us with first order initial value problems. Suppose we had a linear initial value problem of the nth order:
such that
For any nonzero , if and are continuous on some interval containing , exists and is unique.
Related concepts
A delay differential equation (DDE) is an equation for a function of a single variable, usually called time, in which the derivative of the function at a certain time is given in terms of the values of the function at earlier times.
Integral equations may be viewed as the analog to differential equations where instead of the equation involving derivatives, the equation contains integrals.
An integro-differential equation (IDE) is an equation that combines aspects of a differential equation and an integral equation.
A stochastic differential equation (SDE) is an equation in which the unknown quantity is a stochastic process and the equation involves some known stochastic processes, for example, the Wiener process in the case of diffusion equations.
A stochastic partial differential equation (SPDE) is an equation that generalizes SDEs to include space-time noise processes, with applications in quantum field theory and statistical mechanics.
An ultrametric pseudo-differential equation is an equation which contains p-adic numbers in an ultrametric space. Mathematical models that involve ultrametric pseudo-differential equations use pseudo-differential operators instead of differential operators.
A differential algebraic equation (DAE) is a differential equation comprising differential and algebraic terms, given in implicit form.
Connection to difference equations
The theory of differential equations is closely related to the theory of difference equations, in which the coordinates assume only discrete values, and the relationship involves values of the unknown function or functions and values at nearby coordinates. Many methods to compute numerical solutions of differential equations or study the properties of differential equations involve the approximation of the solution of a differential equation by the solution of a corresponding difference equation.
Applications
The study of differential equations is a wide field in pure and applied mathematics, physics, and engineering. All of these disciplines are concerned with the properties of differential equations of various types. Pure mathematics focuses on the existence and uniqueness of solutions, while applied mathematics emphasizes the rigorous justification of the methods for approximating solutions. Differential equations play an important role in modeling virtually every physical, technical, or biological process, from celestial motion, to bridge design, to interactions between neurons. Differential equations such as those used to solve real-life problems may not necessarily be directly solvable, i.e. do not have closed form solutions. Instead, solutions can be approximated using numerical methods.
Many fundamental laws of physics and chemistry can be formulated as differential equations. In biology and economics, differential equations are used to model the behavior of complex systems. The mathematical theory of differential equations first developed together with the sciences where the equations had originated and where the results found application. However, diverse problems, sometimes originating in quite distinct scientific fields, may give rise to identical differential equations. Whenever this happens, mathematical theory behind the equations can be viewed as a unifying principle behind diverse phenomena. As an example, consider the propagation of light and sound in the atmosphere, and of waves on the surface of a pond. All of them may be described by the same second-order partial differential equation, the wave equation, which allows us to think of light and sound as forms of waves, much like familiar waves in the water. Conduction of heat, the theory of which was developed by Joseph Fourier, is governed by another second-order partial differential equation, the heat equation. It turns out that many diffusion processes, while seemingly different, are described by the same equation; the Black–Scholes equation in finance is, for instance, related to the heat equation.
The number of differential equations that have received a name, in various scientific areas is a witness of the importance of the topic. See List of named differential equations.
Software
Some CAS software can solve differential equations. These are the commands used in the leading programs:
Maple: dsolve
Mathematica: DSolve[]
Maxima: ode2(equation, y, x)
SageMath: desolve()
SymPy: sympy.solvers.ode.dsolve(equation)
Xcas: desolve(y'=k*y,y)
| Mathematics | Calculus and analysis | null |
1424877 | https://en.wikipedia.org/wiki/Scotia%20plate | Scotia plate | The Scotia plate () is a minor tectonic plate on the edge of the South Atlantic and Southern oceans. Thought to have formed during the early Eocene with the opening of the Drake Passage that separates Antarctica and South America, it is a minor plate whose movement is largely controlled by the two major plates that surround it: the Antarctic plate and the South American plate. The Scotia plate takes its name from the steam yacht Scotia of the Scottish National Antarctic Expedition (1902–04), the expedition that made the first bathymetric study of the region.
Roughly rhomboid, extending between and , the plate is wide and long. It is moving WSW at /year and the South Sandwich plate is moving east at /year in an absolute reference frame. Its boundaries are defined by the East Scotia Ridge, the North Scotia Ridge, the South Scotia Ridge, and the Shackleton fracture zone.
The Scotia plate is made of oceanic crust and continental fragments now distributed around the Scotia Sea. Before the formation of the plate began (40Ma), these fragments formed a continuous landmass from Patagonia to the Antarctic Peninsula along an active subduction margin. At present, the plate is almost completely submerged, with only the small exceptions of South Georgia on its northeastern edge and the southern tip of South America.
Tectonic setting
Together with the South Sandwich plate, the Scotia plate joins the southernmost Andes to the Antarctic Peninsula, just like the Caribbean plate joins the northernmost Andes to North America, and these two plates are comparable in several ways. Both have volcanic arcs at their eastern ends, the South Sandwich Islands on the Sandwich Plate and the Lesser Antilles on the Caribbean plate, and both plates also had a major impact on global climate when they closed the two major gateways between the Pacific and Atlantic Oceans during the Mesozoic and Cenozoic.
North Scotia Ridge
The northern edge of the Scotia plate is bounded by the South American plate, forming the North Scotia Ridge. The North Scotia ridge is a left-lateral, or sinistral, transform boundary with a transform rate of roughly 7.1 mm/yr. The Magallanes–Fagnano Fault is passing through Tierra del Fuego.
The northern ridge stretches from Isla de los Estados off Tierra del Fuego in the west to the microcontinent South Georgia in the east, with a series of shallow banks in between: Burdwood, Davis, Barker, and Shag Rocks. North of the ridge is the deep Falkland Trough.
South Georgia microcontinent
Experts in plate tectonics have been unable to determine whether the South Georgian Islands are part of the Scotia plate or have been recently accreted to the South American plate. Surface expressions of the plate boundary are found north of the islands suggesting a long-term presence of the transform fault there. Yet seismic studies have identified strain and thrusting south of the islands indicating the possible shift of the transform fault to an area south of the island. It has also been suggested that the plate bearing the islands may have broken off from the Scotia plate, forming a new independent South Georgia microplate, yet there is little evidence to make this conclusion.
The South Georgia microcontinent was originally connected to the Roca Verdes back-arc basin (southernmost Tierra del Fuego) until the Eocene. Before that, this basin went through a series of geological transformations during the Cretaceous, through which South Georgia was first buried, then made a topographic feature again by the Late Cretaceous. At about 45 Ma, South Georgia, still part of the South American plate, got buried again and something, possibly rotation of the Fuegian Andes, completed the break-up and allowed South Georgia a second exhumation. During the Oligocene (34–23 Ma) South Georgia was reburied again as seafloor spreading took place in the West Scotia Sea. 10 Ma, finally, the South Georgia microcontinent was uplifted as a result of the collision with the Northeast Georgia Rise.
South Scotia Ridge
The southern edge of the plate is bordered by the Antarctic plate, forming the South Scotia Ridge, a left-lateral transform boundary sliding at a rate of roughly 7.4–9.5 mm/yr that occupies the southern half of the Antarctic-Scotia plate boundary. The relative motion between the Scotia plate and the Antarctic plate on the western boundary is 7.5–8.7 mm/yr. Though the South Scotia Ridge is overall a transform fault, small sections of the ridge are spreading to make up for the somewhat jagged shape of the boundary.
At the eastern tip of the Antarctic Peninsula, the beginning of South Scotia Ridge, a small and heavily dissected bank degenerates into several outcrops dominated by Paleozoic and Cretaceous rocks. A small basin, Powel Basin, separates this cluster from the South Orkney microcontinent composed of Triassic and younger rocks.
The eastern continuation of the ridge, the Scotia Arc east of the South Sandwich plate, are the South Sandwich island arc and trench. This volcanic active island arc has submerged ancestors in Jane and Discovery banks in the southern ridge.
Shackleton fracture zone
The western edge of the plate is bounded by the Antarctic plate, forming the Shackleton fracture zone and the southern Chile Trench. The Southern Chile Trench is a southern extension of the subduction of the Antarctic and Nazca plates below South America. Heading south along the ridge, the subduction rate decreases until its remaining oblique motion evolves into the Shackleton fracture zone transform boundary. The south-western edge of the plate is bounded by the Shetland microplate separating the Shackleton fracture zone and the South Scotia Ridge.
North of South Shetland Islands and along the southern half of the Shackleton fracture zone is the remnant of the Phoenix plate (also known as Drake or Aluk Plate). Around 47 Ma the subduction of the Phoenix plate started as the propagation of the Pacific-Antarctic Ridge continued. The last collision between Phoenix ridge segments and the subduction zone was 6.5 Ma and at 3.3 Ma movements had stopped and the remnants of the Phoenix plate was incorporated into the Antarctic plate. The southern part of the Shackleton fracture zone is the former eastern edge of the Phoenix plate.
East Scotia Ridge
The eastern edge of the Scotia plate is a spreading ridge bounded by the South Sandwich microplate, forming the East Scotia Ridge. The East Scotia Ridge is a back-arc spreading ridge that formed due to subduction of the South American plate below the South Sandwich plate along the South Sandwich Island arc. Exact spreading rates are still being disputed in the literature, but it has been agreed that rates range between 60 and 90 mm/yr.
The banks of northern Central Scotia Sea are superposed on oceanic basement and the spreading centre of the West Scotia Sea. Analyses of samples of volcaniclastic rocks from these sites indicate they are constructs of a continental arc and in some cases oceanic arc similar to those being formed in the currently active South Sandwich Arc. The oldest volcanic arc activity in the central and eastern regions of the Scotia Sea are 28.5 Ma. The South Sandwich forarc originated in the Central Scotia Sea at that time but has since been translated eastward by the back-arc spreading centre of the East Scotia Ridge.
West Scotia Ridge
One of the most prominent features in the Scotia plate itself is the median valley known as the West Scotia Ridge. It was produced by two plates that are no longer independently active: the Magallanes and Central Scotia plates so is an extinct spreading center. The ridge consists of seven segments separated by right-lateral transform faults of various lengths. The western region of the Scotia plate can be dated to 26–5.5 Ma suggesting spreading
was active before spreading completely translated to the East Scotia Ridge at about 6 Ma. It was assumed that spreading cessation occurred where the W7 segment joins the North Scotia Ridge. However basalt lava sampled from the eastern side of the W7 segment were found to be between 93 and 137 million years old and the geochemistry of arc-like basalts suggests no sea floor spreading actually occurred in the W7 segment itself. This has resulted in the suggestion that the W7 segment is a downfaulted block of the North Scotia Ridge of the Fuegian Andes continental margin arc, or is potentially related to the putative Cretaceous Central Scotia Sea.
Timeline
The timing of the formation of the Scotia plate and opening of the Drake Passage have long been the subject of much debate due to the important implications for changes in ocean currents and shifts in paleoclimate. The thermal isolation of Antarctica, engendering the formation of the Antarctic ice sheet, has largely been attributed to the opening of the Drake Passage.
Formation
The Scotia plate originated about (Ma), during the late Mesozoic at the Panthalassic margins of the Gondwana supercontinent between two Precambrian cratons, the Kalahari and East Antarctic Cratons, now located in Africa and Antarctica. Its development was also influenced by the Río de la Plata Craton in South America. The initial cause for its formation was the break-up of Gondwana in what became the south-west Indian Ocean.
The earliest marine fossils found on the Maurice Ewing Bank, on the eastern end of the Falkland Plateau, are associated with the Indian and Tethys Oceans and probably slightly more than 150 Ma. The Weddell Sea opened and spread along the southern margin of the Falkland Plateau and into the Rocas Verdes back-arc basin which extends from South Georgia and along the Patagonian Andes. Fragments of the inverted oceanic basement of this basin are preserved as ophiolitic complexes in this area, including the Larsen Harbour complex on South Georgia. This makes it possible to restore the original position of the South Georgia microcontinent south of the Burdwood Bank on the western North Scotia Ridge south of the Falkland Islands.
The Rocas Verdes basin was filled with turbidites derived from the volcanic arc on its Pacific margin and partly from its continental margin. The Weddell Sea continued to expand which led to the extension between the Patagonian Andes and the Antarctic Peninsula. In the Mid-Cretaceous (100 Ma) the spreading rate in South Atlantic increased significantly and the Mid-Atlantic Ridge grew from 1200 km to 7000 km. This led to compressional deformations along the western margin of the South American plate and the obduction of the Rocas Verdes basement onto this margin. Structures associated with this obduction are found from Tierra del Fuego to South Georgia. The acceleration of the westward motion of South America and the inversion of the Rocas Verde Basin finally lead to the initiation of the Scotia Arc. This inversion had a strike-slip component which can be seen in the Cooper Bay dislocation on South Georgia. This geological regime lead to the uplift and elongation of the Andes and the embryonic North Scotia Ridge, which resulted in the initial eastward relocation of the South Georgia microcontinent and formation of the Central Scotia Sea.
Opening of Drake Passage
Between the Late Cretaceous and the Early Oligocene (90–30 Ma) little changed in the region, except for the subduction of the Phoenix plate along the Shackleton fracture zone. The Late Paleocene and Early Eocene (60–50 Ma) saw the formation of South Scotia Sea and South Scotia Ridge — the first sign of separation of the southern Andes and the Antarctic Peninsula — which resulted in seafloor spreading in the West Scotia Sea and hence the initial opening of a deep Drake Passage. The ongoing lengthening of the North Scotia Ridge beyond the Burdwood Bank caused South Georgia to move further east. The banks of the North Scotia Ridge contain volcanic rocks similar to those found in Tierra del Fuego, including on Isla de los Estados on the easternmost tip.
During the early Eocene (50 Ma), the Drake Passage between the southern tip of South America at Cape Horn and the South Shetland Islands of Antarctica was a small opening with limited circulation. A change in relative motion between the South American plate and the Antarctic plate would have severe effects, causing seafloor spreading and the formation of the Scotia plate. Marine geophysical data indicates that motion between the South American plate and the Antarctic plate shifted from N–S to WNW–ESE accompanied by an eightfold increase in the separation rate. This shift in spreading initiated crustal thinning and by 30–34 Ma, the West Scotia Ridge formed.
A complex pattern of spreading prior to 26 Ma is probably present in the oldest parts of this spreading regime; i.e. west of Terror Rise (north of Elephant Island) and on the shelf slope of Tierra del Fuego. Until 17 Ma the Central Scotia plate moved quickly eastwards, fuelled by the eastern trench migration, but both plates have moved very slowly since.
| Physical sciences | Tectonic plates | Earth science |
1425422 | https://en.wikipedia.org/wiki/Electric%20arc%20furnace | Electric arc furnace | An electric arc furnace (EAF) is a furnace that heats material by means of an electric arc.
Industrial arc furnaces range in size from small units of approximately one-tonne capacity (used in foundries for producing cast iron products) up to about 400-tonne units used for secondary steelmaking. Arc furnaces used in research laboratories and by dentists may have a capacity of only a few dozen grams. Industrial electric arc furnace temperatures can reach , while laboratory units can exceed .
In electric arc furnaces, the charge material (the material entered into the furnace for heating, not to be confused with electric charge) is directly exposed to an electric arc, and the current from the electrode terminals passes through the charge material.
Arc furnaces differ from induction furnaces, in which the charge is heated instead by eddy currents.
History
In the 19th century, a number of people had employed an electric arc to melt iron. Sir Humphry Davy conducted an experimental demonstration in 1810; welding was investigated by Pepys in 1815; Pinchon attempted to create an electrothermic furnace in 1853; and, in 1878–79, Sir William Siemens took out patents for electric furnaces of the arc type.
The first successful and operational furnace was invented by James Burgess Readman in Edinburgh, Scotland, in 1888 and patented in 1889. This was specifically for the creation of phosphorus.
Further electric arc furnaces were developed by Paul Héroult, of France, with a commercial plant established in the United States in 1907. The Sanderson brothers formed The Sanderson Brothers Steel Co. in Syracuse, New York, installing the first electric arc furnace in the U.S. This furnace is now on display at Station Square, Pittsburgh, Pennsylvania.
Initially "electric steel" produced by an electric arc furnace was a specialty product for such uses as machine tools and spring steel. Arc furnaces were also used to prepare calcium carbide for use in carbide lamps. The Stassano electric furnace is an arc type furnace that usually rotates to mix the bath. The Girod furnace is similar to the Héroult furnace.
While EAFs were widely used in World War II for production of alloy steels, it was only later that electric steelmaking began to expand. The low capital cost for a mini-mill—around US$140–200 per ton of annual installed capacity, compared with US$1,000 per ton of annual installed capacity for an integrated steel mill—allowed mills to be quickly established in war-ravaged Europe, and also allowed them to successfully compete with the big United States steelmakers, such as Bethlehem Steel and U.S. Steel, for low-cost, carbon steel "long products" (structural steel, rod and bar, wire, and fasteners) in the U.S. market.
When Nucor—now one of the largest steel producers in the US — entered the market for long steel products in 1969, they used a mini-mill with an EAF as its steelmaking furnace, soon followed by other manufacturers. While Nucor expanded rapidly in the Eastern US, the companies that followed them into mini-mill operations concentrated on local markets for long products, where the EAF allowed the plants to vary production according to local demand. This pattern was followed globally, with EAF steel production primarily used for long products, while integrated mills, using blast furnaces and basic oxygen furnaces, cornered the markets for "flat products"—sheet steel and heavier steel plate. In 1987, Nucor expanded into the flat products market, still using the EAF production method.
Construction
An electric arc furnace used for steelmaking consists of a refractory-lined vessel, usually water-cooled in larger sizes, covered with a retractable roof, and through which one or more graphite electrodes enter the furnace.
The furnace is primarily split into three sections:
the shell, which consists of the sidewalls and lower steel "bowl";
the hearth, which consists of the refractory that lines the lower bowl;
the roof, which may be refractory-lined or water-cooled, and can be shaped as a section of a sphere, or as a frustum (conical section). The roof also supports the refractory delta in its centre, through which one or more graphite electrodes enter.
The hearth may be hemispherical in shape, or in an eccentric bottom tapping furnace (see below), the hearth has the shape of a halved egg. In modern meltshops, the furnace is often raised off the ground floor, so that ladles and slag pots can easily be maneuvered under either end of the furnace. Separate from the furnace structure is the electrode support and electrical system, and the tilting platform on which the furnace rests. Two configurations are possible: the electrode supports and the roof tilt with the furnace, or are fixed to the raised platform.
A typical alternating current furnace is powered by a three-phase electrical supply, and therefore has three electrodes. Electrodes are round in section, and typically in segments with threaded couplings, so that as the electrodes wear, new segments can be added. The arc forms between the charged material and the electrode; the charge is heated both by current passing through the charge and by the radiant energy evolved by the arc. The electric arc temperature reaches around , thus causing the lower sections of the electrodes to glow incandescently when in operation. The electrodes are automatically raised and lowered by a positioning system, which may use either electric winch hoists or hydraulic cylinders. The regulating system maintains approximately constant current and power input during the melting of the charge, even though scrap may move under the electrodes as it melts. The mast arms holding the electrodes can either carry heavy busbars (which may be hollow water-cooled copper pipes carrying current to the electrode clamps) or be "hot arms", where the whole arm carries the current, increasing efficiency. Hot arms can be made from copper-clad steel or aluminium. Large water-cooled cables connect the bus tubes or arms with the transformer located adjacent to the furnace. The transformer is installed in a vault and is cooled by pump-circulated transformer oil, with the oil being cooled by water via heat exchangers.
The furnace is built on a tilting platform so that the liquid steel can be poured into another vessel for transport. The operation of tilting the furnace to pour molten steel is called "tapping". Originally, all steelmaking furnaces had a tapping spout closed with refractory that washed out when the furnace was tilted, but often modern furnaces have an eccentric bottom tap-hole (EBT) to reduce inclusion of nitrogen and slag in the liquid steel. These furnaces have a taphole that passes vertically through the hearth and shell, and is set off-centre in the narrow "nose" of the egg-shaped hearth. It is filled with refractory sand, such as olivine, when it is closed off. Modern plants may have two shells with a single set of electrodes that can be transferred between the two; one shell preheats scrap while the other shell is utilised for meltdown. Other DC-based furnaces have a similar arrangement, but have electrodes for each shell and one set of electronics.
AC furnaces usually exhibit a pattern of hot and cold-spots around the hearth perimeter, with the cold-spots located between the electrodes. Modern furnaces mount oxygen-fuel burners in the sidewall and use them to provide chemical energy to the cold-spots, making the heating of the steel more uniform. Additional chemical energy is provided by injecting oxygen and carbon into the furnace; historically this was done through lances (hollow mild-steel tubes) in the slag door, but now this is mainly done through wall-mounted injection units that combine the oxygen-fuel burners and the oxygen or carbon injection systems into one unit.
A mid-sized modern steelmaking furnace would have a transformer rated about 60,000,000 volt-amperes (60 MVA), with a secondary voltage between 400 and 900 volts and a secondary current in excess of 44,000 amperes. In a modern shop such a furnace would be expected to produce a quantity of 80 tonnes of liquid steel in approximately 50 minutes from charging with cold scrap to tapping the furnace. In comparison, basic oxygen furnaces can have a capacity of 150–300 tonnes per batch, or "heat", and can produce a heat in 30–40 minutes. Enormous variations exist in furnace design details and operation, depending on the end product and local conditions, as well as ongoing research to improve furnace efficiency. The largest scrap-only furnace (in terms of tapping weight and transformer rating) is a DC furnace operated by Tokyo Steel in Japan, with a tap weight of 420 tonnes and fed by eight 32 MVA transformers for 256 MVA total power.
Energy density
To produce a ton of steel in an electric arc furnace requires approximately 400 kilowatt-hours (1.44 gigajoules) per short ton or about 440 kWh (1.6 GJ) per tonne. The theoretical minimum amount of energy required to melt a tonne of scrap steel is 300 kWh (1.09 GJ) (melting point ). Therefore, a 300-tonne, 300 MVA EAF will require approximately 132 MWh of energy to melt the steel, and a "power-on time" (the time that steel is being melted with an arc) of approximately 37 minutes.
Electric arc steelmaking is only economical where there is plentiful, reliable electricity, with a well-developed electrical grid. In many locations, mills operate during off-peak hours when utilities have surplus power generating capacity and the price of electricity is less. This compares very favourably with energy consumption of global steel production by all methods estimated at some 5,555 kWh (20 GJ) per tonne (1 gigajoule is equal to approximately 270 kWh).
Operation
Scrap metal is delivered to a scrap bay, located next to the melt shop. Scrap generally comes in two main grades: shred (whitegoods, cars and other objects made of similar light-gauge steel) and heavy melt (large slabs and beams), along with some direct reduced iron (DRI) or pig iron for chemical balance. Some furnaces melt almost 100% DRI.
The scrap is loaded into large buckets called baskets, with "clamshell" doors for a base. Care is taken to layer the scrap in the basket to ensure good furnace operation; heavy melt is placed on top of a light layer of protective shred, on top of which is placed more shred. These layers should be present in the furnace after charging. After loading, the basket may pass to a scrap pre-heater, which uses hot furnace off-gases to heat the scrap and recover energy, increasing plant efficiency.
The scrap basket is then taken to the melt shop, the roof is swung off the furnace, and the furnace is charged with scrap from the basket. Charging is one of the more dangerous operations for the EAF operators. A lot of potential energy is released by the tonnes of falling metal; any liquid metal in the furnace is often displaced upwards and outwards by the solid scrap, and the grease and dust on the scrap is ignited if the furnace is hot, resulting in a fireball erupting.
In some twin-shell furnaces, the scrap is charged into the second shell while the first is being melted down, and pre-heated with off-gas from the active shell. Other operations are continuous charging—pre-heating scrap on a conveyor belt, which then discharges the scrap into the furnace proper, or charging the scrap from a shaft set above the furnace, with off-gases directed through the shaft. Other furnaces can be charged with hot (molten) metal from other operations.
After charging, the roof is swung back over the furnace and meltdown commences. The electrodes are lowered onto the scrap, an arc is struck and the electrodes are then set to bore into the layer of shred at the top of the furnace. Lower voltages are selected for this first part of the operation to protect the roof and walls from excessive heat and damage from the arcs. Once the electrodes have reached the heavy melt at the base of the furnace and the arcs are shielded by the scrap, the voltage can be increased and the electrodes raised slightly, lengthening the arcs and increasing power to the melt. This enables a molten pool to form more rapidly, reducing tap-to-tap times. Oxygen is blown into the scrap, combusting or cutting the steel, and extra chemical heat is provided by wall-mounted oxygen-fuel burners. Both processes accelerate scrap meltdown. Supersonic nozzles enable oxygen jets to penetrate foaming slag and reach the liquid bath.
An important part of steelmaking is the formation of slag, which floats on the surface of the molten steel. Slag usually consists of metal oxides, and acts as a destination for oxidised impurities, as a thermal blanket (stopping excessive heat loss) and helping to reduce erosion of the refractory lining. For a furnace with basic refractories, which includes most carbon steel-producing furnaces, the usual slag formers are calcium oxide (CaO, in the form of burnt lime) and magnesium oxide (MgO, in the form of dolomite and magnesite).
These slag formers are either charged with the scrap, or blown into the furnace during meltdown. Another major component of EAF slag is iron oxide from steel combusting with the injected oxygen. Later in the heat, carbon (in the form of coke or coal) is injected into this slag layer, reacting with the iron oxide to form metallic iron and carbon monoxide gas, which then causes the slag to foam, allowing greater thermal efficiency, and better arc stability and electrical efficiency. The slag blanket also covers the arcs, preventing damage to the furnace roof and sidewalls from radiant heat.
Once the initial scrap charge has been melted down, another bucket of scrap can be charged into the furnace, although EAF development is moving towards single-charge designs. The scrap-charging and meltdown process can be repeated as many times as necessary to reach the required heat weight - the number of charges is dependent on the density of scrap; lower-density scrap means more charges. After all scrap charges have completely melted, refining operations take place to check and correct the steel chemistry and superheat the melt above its freezing temperature in preparation for tapping.
More slag formers are introduced and more oxygen is blown into the bath, burning out impurities such as silicon, sulfur, phosphorus, aluminium, manganese, and calcium, and removing their oxides to the slag. Removal of carbon takes place after these elements have burnt out first, as they have a greater affinity for oxygen. Metals that have a poorer affinity for oxygen than iron, such as nickel and copper, cannot be removed through oxidation and must be controlled through scrap chemistry alone, such as introducing the direct reduced iron and pig iron mentioned earlier.
A foaming slag is maintained throughout, and often overflows the furnace to pour out of the slag door into the slag pit. Temperature sampling and chemical sampling take place via automatic lances. Oxygen and carbon can be automatically measured via special probes that dip into the steel, but for all other elements, a "chill" sample — a small, solidified sample of the steel — is analysed on an arc-emission spectrometer.
Once the temperature and chemistry are correct, the steel is tapped out into a preheated ladle through tilting the furnace. For plain-carbon steel furnaces, as soon as slag is detected during tapping the furnace is rapidly tilted back towards the deslagging side, minimising slag carryover into the ladle. For some special steel grades, including stainless steel, the slag is poured into the ladle as well, to be treated at the ladle furnace to recover valuable alloying elements. During tapping some alloy additions are introduced into the metal stream, and more fluxes such as lime are added on top of the ladle to begin building a new slag layer.
Often, a few tonnes of liquid steel and slag is left in the furnace in order to form a "hot heel", which helps preheat the next charge of scrap and accelerate its meltdown. During and after tapping, the furnace is "turned around": the slag door is cleaned of solidified slag, the visible refractories are inspected and water-cooled components checked for leaks, and electrodes are inspected for damage or lengthened through the addition of new segments. The taphole is filled with sand at the completion of tapping. For a 90-tonne, medium-power furnace, the whole process will usually take about 60–70 minutes from the tapping of one heat to the tapping of the next (the tap-to-tap time).
The furnace is completely emptied of steel and slag on a regular basis so that an inspection of the refractories can be made and larger repairs made if necessary. As the refractories are often made from calcined carbonates, they are extremely susceptible to hydration from water, so any suspected leaks from water-cooled components are treated extremely seriously, beyond the immediate concern of potential steam explosions. Excessive refractory wear can lead to breakouts, where the liquid metal and slag penetrate the refractory and furnace shell and escape into the surrounding areas.
Advantages for steelmaking
The use of EAFs allows steel to be made from a 100% scrap metal feedstock. This greatly reduces the energy required to make steel when compared with primary steelmaking from ores.
Another benefit is flexibility: while blast furnaces cannot vary their production by much and can remain in operation for years at a time, EAFs can be rapidly started and stopped, allowing the steel mill to vary production according to demand.
Although steelmaking arc furnaces generally use scrap steel as their primary feedstock, if hot metal from a blast furnace or direct-reduced iron is available economically, these can also be used as furnace feed.
As EAFs require large amounts of electrical power, many companies schedule their operations to take advantage of off-peak electricity pricing.
A typical steelmaking arc furnace is the source of steel for a mini-mill, which may make bars or strip product. Mini-mills can be sited relatively near the markets for steel products, so the transport requirements are less than for an integrated mill, which would commonly be sited near a harbor for better access to shipping.
Depending on the proportions of steel scrap, DRI and pig iron used, electric arc furnace steelmaking can result in carbon dioxide emissions as low as 0.6 tons CO2 per ton of steel produced, which is significantly lower than the conventional production route via blast furnaces and the basic oxygen furnace, which produces 2.9 tons CO2 per ton of steel produced.
Issues
Although the modern electric arc furnace is a highly efficient recycler of steel scrap, operation of an arc furnace shop can have adverse environmental effects. Much of the capital cost of a new installation will be devoted to systems intended to reduce these effects, which include:
Enclosures to reduce high sound levels
Dust collector for furnace off-gas
slag production
cooling water demand
Heavy truck traffic for scrap, materials handling, and product
Environmental effects of electricity generation
Since EAF steelmaking mainly use recycled materials like scrap iron and scrap steel, as their composition varies the resulting EAF slag and EAF dust can be toxic. EAF dust is collected by air pollution control equipment. It is called collected dust and usually contains heavy metals, such as zinc, lead and dioxins, etc. It is categorized as hazardous industrial waste and disposal is regulated.
Because of the very dynamic quality of the arc furnace load, power systems may require technical measures to maintain the quality of power for other customers; flicker and harmonic distortion are common power system side-effects of arc furnace operation.
Other electric arc furnaces
For steelmaking, direct current (DC) arc furnaces are used, with a single electrode in the roof and the current return through a conductive bottom lining or conductive pins in the base. The advantage of DC is lower electrode consumption per ton of steel produced, since only one electrode is used, as well as less electrical harmonics and other similar problems. The size of DC arc furnaces is limited by the current carrying capacity of available electrodes, and the maximum allowable voltage. Maintenance of the conductive furnace hearth is a bottleneck in extended operation of a DC arc furnace.
In a steel plant, a ladle furnace (LF) is used to maintain the temperature of liquid steel during processing after tapping from EAF or to change the alloy composition. The ladle is used for the first purpose when there is a delay later in the steelmaking process. The ladle furnace consists of a refractory roof, a heating system, and, when applicable, a provision for injecting argon gas into the bottom of the melt for stirring. Unlike a scrap melting furnace, a ladle furnace does not have a tilting or scrap-charging mechanism.
Electric arc furnaces are also used for production of calcium carbide, ferroalloys, and other non-ferrous alloys, and for production of phosphorus. Furnaces for these services are physically different from steel-making furnaces and may operate on a continuous, rather than batch, basis. Continuous-process furnaces may also use paste-type, Søderberg electrodes to prevent interruptions from electrode changes.
Such a furnace is known as a submerged arc furnace, because the electrode tips are buried in the slag/charge, and arcing occurs through the slag, between the matte and the electrode. The casing and casing fins of the electrode melt the electrode paste through electrical current passing through the electrode casing and heat from the furnace. A steelmaking arc furnace, by comparison, arcs in the open. The key is the electrical resistance, which is what generates the heat required: the resistance in a steelmaking furnace is the atmosphere, while in a submerged-arc furnace, the slag (or charge) supplies the resistance. The liquid metal formed in either furnace is too conductive to form an effective heat-generating resistance.
Amateurs have constructed a variety of arc furnaces, often based on electric arc welding kits contained by silica blocks or flower pots. Though crude, these simple furnaces can melt a wide range of materials, create calcium carbide, and more.
Cooling methods
Smaller arc furnaces may be adequately cooled by circulation of air over structural elements of the shell and roof, but larger installations require intensive forced cooling to maintain the structure within safe operating limits. The furnace shell and roof may be cooled either by water circulated through pipes which form a panel, or by water sprayed on the panel elements. Tubular panels may be replaced when they become cracked or reach their thermal stress life cycle.
Spray cooling is the most economical and is the highest efficiency cooling method. A spray cooling piece of equipment can be relined almost endlessly. Equipment that lasts 20 years is the norm. While a tubular leak is immediately noticed in an operating furnace due to the pressure loss alarms on the panels, at this time there exists no immediate way of detecting a very small volume spray cooling leak. These typically hide behind slag coverage and can hydrate the refractory in the hearth, leading to a break out of molten metal or in the worst case a steam explosion.
Plasma arc furnace
A plasma arc furnace (PAF) uses plasma torches instead of graphite electrodes. Each of these torches has a casing with a nozzle and axial tubing for feeding a plasma-forming gas (either nitrogen or argon) and a burnable cylindrical graphite electrode within the tubing. Such furnaces can be called plasma arc melt (PAM) furnaces; they are used extensively in the titanium-melting industry and similar specialty metal industries.
Vacuum arc remelting
Vacuum arc remelting (VAR) is a secondary remelting process for vacuum refining and manufacturing of ingots with improved chemical and mechanical homogeneity.
In critical military and commercial aerospace applications, material engineers commonly specify VIM-VAR steels. VIM means vacuum induction melted and VAR means vacuum arc remelted. VIM-VAR steels become bearings for jet engines, rotor shafts for military helicopters, flap actuators for fighter jets, gears in jet or helicopter transmissions, mounts or fasteners for jet engines, jet tail hooks and other demanding applications.
Most grades of steel are melted once and are then cast or teemed into a solid form prior to extensive forging or rolling to a metallurgically-sound form. In contrast, VIM-VAR steels go through two more highly purifying melts under vacuum. After melting in an electric arc furnace and alloying in an argon oxygen decarburization vessel, steels destined for vacuum remelting are cast into ingot molds. The solidified ingots then head for a vacuum induction melting furnace. This vacuum remelting process rids the steel of inclusions and unwanted gases while optimizing the chemical composition.
The VIM operation returns these solid ingots to the molten state in the contaminant-free void of a vacuum. This tightly controlled melt often requires up to 24 hours. Still enveloped by the vacuum, the hot metal flows from the VIM furnace crucible into giant electrode molds. A typical electrode is about 15 feet (5 m) tall and will be in various diameters. The electrodes solidify under vacuum.
For VIM-VAR steels, the surface of the cooled electrodes must be ground to remove surface irregularities and impurities before the next vacuum remelt. Then the ground electrode is placed in a VAR furnace. In a VAR furnace, the steel gradually melts drop-by-drop in the vacuum-sealed chamber. Vacuum arc remelting further removes lingering inclusions to provide superior steel cleanliness and remove gases like oxygen, nitrogen and hydrogen. Controlling the rate at which these droplets form and solidify ensures a consistency of chemistry and microstructure throughout the entire VIM-VAR ingot, making the steel more resistant to fracture or fatigue. This refinement process is essential to meet the performance characteristics of parts like a helicopter rotor shaft, a flap actuator on a military jet, or a bearing in a jet engine.
For some commercial or military applications, steel alloys may go through only one vacuum remelt, namely the VAR. For example, steels for solid rocket cases, landing gears, or torsion bars for fighting vehicles typically involve one vacuum remelt.
Vacuum arc remelting is also used in production of titanium and other metals which are reactive or in which high purity is required.
| Technology | Metallurgy | null |
1425534 | https://en.wikipedia.org/wiki/Multipole%20expansion | Multipole expansion | A multipole expansion is a mathematical series representing a function that depends on angles—usually the two angles used in the spherical coordinate system (the polar and azimuthal angles) for three-dimensional Euclidean space, . Multipole expansions are useful because, similar to Taylor series, oftentimes only the first few terms are needed to provide a good approximation of the original function. The function being expanded may be real- or complex-valued and is defined either on , or less often on for some other
Multipole expansions are used frequently in the study of electromagnetic and gravitational fields, where the fields at distant points are given in terms of sources in a small region. The multipole expansion with angles is often combined with an expansion in radius. Such a combination gives an expansion describing a function throughout three-dimensional space.
The multipole expansion is expressed as a sum of terms with progressively finer angular features (moments). The first (the zeroth-order) term is called the monopole moment, the second (the first-order) term is called the dipole moment, the third (the second-order) the quadrupole moment, the fourth (third-order) term is called the octupole moment, and so on. Given the limitation of Greek numeral prefixes, terms of higher order are conventionally named by adding "-pole" to the number of poles—e.g., 32-pole (rarely dotriacontapole or triacontadipole) and 64-pole (rarely tetrahexacontapole or hexacontatetrapole). A multipole moment usually involves powers (or inverse powers) of the distance to the origin, as well as some angular dependence.
In principle, a multipole expansion provides an exact description of the potential, and generally converges under two conditions: (1) if the sources (e.g. charges) are localized close to the origin and the point at which the potential is observed is far from the origin; or (2) the reverse, i.e., if the sources are located far from the origin and the potential is observed close to the origin. In the first (more common) case, the coefficients of the series expansion are called exterior multipole moments or simply multipole moments whereas, in the second case, they are called interior multipole moments.
Expansion in spherical harmonics
Most commonly, the series is written as a sum of spherical harmonics. Thus, we might write a function as the sum
where are the standard spherical harmonics, and are constant coefficients which depend on the function. The term represents the monopole; represent the dipole; and so on. Equivalently, the series is also frequently written as
where the represent the components of a unit vector in the direction given by the angles and , and indices are implicitly summed. Here, the term is the monopole; is a set of three numbers representing the dipole; and so on.
In the above expansions, the coefficients may be real or complex. If the function being expressed as a multipole expansion is real, however, the coefficients must satisfy certain properties. In the spherical harmonic expansion, we must have
In the multi-vector expansion, each coefficient must be real:
While expansions of scalar functions are by far the most common application of multipole expansions, they may also be generalized to describe tensors of arbitrary rank. This finds use in multipole expansions of the vector potential in electromagnetism, or the metric perturbation in the description of gravitational waves.
For describing functions of three dimensions, away from the coordinate origin, the coefficients of the multipole expansion can be written as functions of the distance to the origin, —most frequently, as a Laurent series in powers of . For example, to describe the electromagnetic potential, , from a source in a small region near the origin, the coefficients may be written as:
Applications
Multipole expansions are widely used in problems involving gravitational fields of systems of masses, electric and magnetic fields of charge and current distributions, and the propagation of electromagnetic waves. A classic example is the calculation of the exterior multipole moments of atomic nuclei from their interaction energies with the interior multipoles of the electronic orbitals. The multipole moments of the nuclei report on the distribution of charges within the nucleus and, thus, on the shape of the nucleus. Truncation of the multipole expansion to its first non-zero term is often useful for theoretical calculations.
Multipole expansions are also useful in numerical simulations, and form the basis of the fast multipole method of Greengard and Rokhlin, a general technique for efficient computation of energies and forces in systems of interacting particles. The basic idea is to decompose the particles into groups; particles within a group interact normally (i.e., by the full potential), whereas the energies and forces between groups of particles are calculated from their multipole moments. The efficiency of the fast multipole method is generally similar to that of Ewald summation, but is superior if the particles are clustered, i.e. the system has large density fluctuations.
Multipole expansion of a potential outside an electrostatic charge distribution
Consider a discrete charge distribution consisting of point charges with position vectors . We assume the charges to be clustered around the origin, so that for all i: , where has some finite value. The potential , due to the charge distribution, at a point outside the charge distribution, i.e., , can be expanded in powers of . Two ways of making this expansion can be found in the literature: The first is a Taylor series in the Cartesian coordinates , , and , while the second is in terms of spherical harmonics which depend on spherical polar coordinates. The Cartesian approach has the advantage that no prior knowledge of Legendre functions, spherical harmonics, etc., is required. Its disadvantage is that the derivations are fairly cumbersome (in fact a large part of it is the implicit rederivation of the Legendre expansion of , which was done once and for all by Legendre in the 1780s). Also it is difficult to give a closed expression for a general term of the multipole expansion—usually only the first few terms are given followed by an ellipsis.
Expansion in Cartesian coordinates
Assume for convenience. The Taylor expansion of around the origin can be written as
with Taylor coefficients
If satisfies the Laplace equation, then by the above expansion we have
and the expansion can be rewritten in terms of the components of a traceless Cartesian second rank tensor:
where is the Kronecker delta and . Removing the trace is common, because it takes the rotationally invariant out of the second rank tensor.
Example
Consider now the following form of :
Then by direct differentiation it follows that
Define a monopole, dipole, and (traceless) quadrupole by, respectively,
and we obtain finally the first few terms of the multipole expansion of the total potential, which is the sum of the Coulomb potentials of the separate charges:
This expansion of the potential of a discrete charge distribution is very similar to the one in real solid harmonics given below. The main difference is that the present one is in terms of linearly dependent quantities, for
Note:
If the charge distribution consists of two charges of opposite sign which are an infinitesimal distance apart, so that , it is easily shown that the dominant term in the expansion is
the electric dipolar potential field.
Spherical form
The potential at a point outside the charge distribution, i.e. , can be expanded by the Laplace expansion:
where is an irregular solid harmonic (defined below as a spherical harmonic function divided by ) and is a regular solid harmonic (a spherical harmonic times ). We define the spherical multipole moment of the charge distribution as follows
Note that a multipole moment is solely determined by the charge distribution (the positions and magnitudes of the N charges).
A spherical harmonic depends on the unit vector . (A unit vector is determined by two spherical polar angles.) Thus, by definition, the irregular solid harmonics can be written as
so that the multipole expansion of the field at the point outside the charge distribution is given by
This expansion is completely general in that it gives a closed form for all terms, not just for the first few. It shows that the spherical multipole moments appear as coefficients in the expansion of the potential.
It is of interest to consider the first few terms in real form, which are the only terms commonly found in undergraduate textbooks.
Since the summand of the m summation is invariant under a unitary transformation of both factors simultaneously and since transformation of complex spherical harmonics to real form is by a unitary transformation, we can simply substitute real irregular solid harmonics and real multipole moments. The term becomes
This is in fact Coulomb's law again. For the term we introduce
Then
This term is identical to the one found in Cartesian form.
In order to write the term, we have to introduce shorthand notations for the five real components of the quadrupole moment and the real spherical harmonics. Notations of the type
can be found in the literature. Clearly the real notation becomes awkward very soon, exhibiting the usefulness of the complex notation.
Interaction of two non-overlapping charge distributions
Consider two sets of point charges, one set clustered around a point and one set clustered around a point . Think for example of two molecules, and recall that a molecule by definition consists of electrons (negative point charges) and nuclei (positive point charges). The total electrostatic interaction energy between the two distributions is
This energy can be expanded in a power series in the inverse distance of and .
This expansion is known as the multipole expansion of UAB.
In order to derive this multipole expansion, we write , which is a vector pointing from towards . Note that
We assume that the two distributions do not overlap:
Under this condition we may apply the Laplace expansion in the following form
where and are irregular and regular solid harmonics, respectively. The translation of the regular solid harmonic gives a finite expansion,
where the quantity between pointed brackets is a Clebsch–Gordan coefficient. Further we used
Use of the definition of spherical multipoles and covering of the summation ranges in a somewhat different order (which is only allowed for an infinite range of ) gives finally
This is the multipole expansion of the interaction energy of two non-overlapping charge distributions which are a distance RAB apart. Since
this expansion is manifestly in powers of . The function is a normalized spherical harmonic.
Molecular moments
All atoms and molecules (except S-state atoms) have one or more non-vanishing permanent multipole moments. Different definitions can be found in the literature, but the following definition in spherical form has the advantage that it is contained in one general equation. Because it is in complex form it has as the further advantage that it is easier to manipulate in calculations than its real counterpart.
We consider a molecule consisting of N particles (electrons and nuclei) with charges eZi. (Electrons have a Z-value of −1, while for nuclei it is the atomic number). Particle i has spherical polar coordinates ri, θi, and φi and Cartesian coordinates xi, yi, and zi.
The (complex) electrostatic multipole operator is
where is a regular solid harmonic function in Racah's normalization (also known as Schmidt's semi-normalization).
If the molecule has total normalized wave function Ψ (depending on the coordinates of electrons and nuclei), then the multipole moment of order of the molecule is given by the expectation (expected) value:
If the molecule has certain point group symmetry, then this is reflected in the wave function: Ψ transforms according to a certain irreducible representation λ of the group ("Ψ has symmetry type λ"). This has the consequence that selection rules hold for the expectation value of the multipole operator, or in other words, that the expectation value may vanish because of symmetry. A well-known example of this is the fact that molecules with an inversion center do not carry a dipole (the expectation values of vanish for . For a molecule without symmetry, no selection rules are operative and such a molecule will have non-vanishing multipoles of any order (it will carry a dipole and simultaneously a quadrupole, octupole, hexadecapole, etc.).
The lowest explicit forms of the regular solid harmonics (with the Condon-Shortley phase) give:
(the total charge of the molecule). The (complex) dipole components are:
Note that by a simple linear combination one can transform the complex multipole operators to real ones. The real multipole operators are of cosine type
or sine type . A few of the lowest ones are:
Note on conventions
The definition of the complex molecular multipole moment given above is the complex conjugate of the definition given in this article, which follows the definition of the standard textbook on classical electrodynamics by Jackson, except for the normalization. Moreover, in the classical definition of Jackson the equivalent of the N-particle quantum mechanical expectation value is an integral over a one-particle charge distribution. Remember that in the case of a one-particle quantum mechanical system the expectation value is nothing but an integral over the charge distribution (modulus of wavefunction squared), so that the definition of this article is a quantum mechanical N-particle generalization of Jackson's definition.
The definition in this article agrees with, among others, the one of Fano and Racah and Brink and Satchler.
Examples
There are many types of multipole moments, since there are many types of potentials and many ways of approximating a potential by a series expansion, depending on the coordinates and the symmetry of the charge distribution. The most common expansions include:
Axial multipole moments of a potential;
Spherical multipole moments of a potential; and
Cylindrical multipole moments of a potential
Examples of potentials include the electric potential, the magnetic potential and the gravitational potential of point sources. An example of a potential is the electric potential of an infinite line charge.
General mathematical properties
Multipole moments in mathematics and mathematical physics form an orthogonal basis for the decomposition of a function, based on the response of a field to point sources that are brought infinitely close to each other. These can be thought of as arranged in various geometrical shapes, or, in the sense of distribution theory, as directional derivatives.
Multipole expansions are related to the underlying rotational symmetry of the physical laws and their associated differential equations. Even though the source terms (such as the masses, charges, or currents) may not be symmetrical, one can expand them in terms of irreducible representations of the rotational symmetry group, which leads to spherical harmonics and related sets of orthogonal functions. One uses the technique of separation of variables to extract the corresponding solutions for the radial dependencies.
In practice, many fields can be well approximated with a finite number of multipole moments (although an infinite number may be required to reconstruct a field exactly). A typical application is to approximate the field of a localized charge distribution by its monopole and dipole terms. Problems solved once for a given order of multipole moment may be linearly combined to create a final approximate solution for a given source.
| Mathematics | Sequences and series | null |
31557892 | https://en.wikipedia.org/wiki/Ductility%20%28Earth%20science%29 | Ductility (Earth science) | In Earth science, ductility refers to the capacity of a rock to deform to large strains without macroscopic fracturing. Such behavior may occur in unlithified or poorly lithified sediments, in weak materials such as halite or at greater depths in all rock types where higher temperatures promote crystal plasticity and higher confining pressures suppress brittle fracture. In addition, when a material is behaving ductilely, it exhibits a linear stress vs strain relationship past the elastic limit.
Ductile deformation is typically characterized by diffuse deformation (i.e. lacking a discrete fault plane) and on a stress-strain plot is accompanied by steady state sliding at failure, compared to the sharp stress drop observed in experiments during brittle failure.
Brittle–Ductile Transition Zone
The brittle–ductile transition zone is characterized by a change in rock failure mode, at an approximate average depth of 10–15 km (~ 6.2–9.3 miles) in continental crust, below which rock becomes less likely to fracture and more likely to deform ductilely. The zone exists because as depth increases confining pressure increases, and brittle strength increases with confining pressure whilst ductile strength decreases with increasing temperature. The transition zone occurs at the point where brittle strength equals ductile strength. In glacial ice this zone is at approximately depth.
Not all materials, however, abide by this transition. It is possible and not rare for material above the transition zone to deform ductilely, and for material below to deform in a brittle manner. The depth of the material does exert an influence on the mode of deformation, but other substances, such as loose soils in the upper crust, malleable rocks, biological debris, and more are just a few examples of that which does not deform in accordance to the transition zone.
The type of dominating deformation process also has a great impact on the types of rocks and structures found at certain depths within the Earth's crust. As evident from Fig. 1.1, different geological formations and rocks are found in accordance to the dominant deformation process. Gouge and Breccia form in the uppermost, brittle regime while Cataclasite and Pseudotachylite form in the lower parts of the brittle regime, edging upon the transition zone. Mylonite forms in the more ductile regime at greater depths while Blastomylonite forms well past the transition zone and well into the ductile regime, even deeper into the crust.
Quantification
Ductility is a material property that can be expressed in a variety of ways. Mathematically, it is commonly expressed as a total quantity of elongation or a total quantity of the change in cross sectional area of a specific rock until macroscopic brittle behavior, such as fracturing, is observed. For accurate measurement, this must be done under several controlled conditions, including but not limited to Pressure, Temperature, Moisture Content, Sample Size, etc., for all can impact the measured ductility. It is important to understand that even the same type of rock or mineral may exhibit different behavior and degrees of ductility due to internal heterogeneities small scale differences between each individual sample. The two quantities are expressed in the form of a ratio or a percent.
% Elongation of a Rock =
Where:
= Initial Length of Rock
= Final Length of Rock
% Change in Area of a Rock =
Where:
= Initial Area
= Final Area
For each of these methods of quantifying, one must take measurements of both the initial and final dimensions of the rock sample. For Elongation, the measurement is a uni-dimensional initial and final length, the former measured before any Stress is applied and the latter measuring the length of the sample after fracture occurs. For Area, it is strongly preferable to use a rock that has been cut into a cylindrical shape before stress application so that the cross-sectional area of the sample can be taken.
Cross-Sectional Area of a Cylinder = Area of a Circle =
Using this, the initial and final areas of the sample can be used to quantify the % change in the area of the rock.
Deformation
Any material is shown to be able to deform ductilely or brittlely, in which the type of deformation is governed by both the external conditions around the rock and the internal conditions sample. External conditions include temperature, confining pressure, presence of fluids, etc. while internal conditions include the arrangement of the crystal lattice, the chemical composition of the rock sample, the grain size of the material, etc.
Ductilely Deformative behavior can be grouped into three categories: Elastic, Viscous, and Crystal-Plastic Deformation.
Elastic Deformation
Elastic Deformation is deformation which exhibits a linear stress-strain relationship (quantified by Young's Modulus) and is derived from Hooke's Law of spring forces (see Fig. 1.2). In elastic deformation, objects show no permanent deformation after the stress has been removed from the system and return to their original state.
Where:
= Stress (In Pascals)
= Young's Modulus (In Pascals)
= Strain (Unitless)
Viscous Deformation
Viscous Deformation is when rocks behave and deform more like a fluid than a solid. This often occurs under great amounts of pressure and at very high temperatures. In viscous deformation, stress is proportional to the strain rate, and each rock sample has its own material property called its Viscosity. Unlike elastic deformation, viscous deformation is permanent even after the stress has been removed.
Where:
= Stress (In Pascals)
= Viscosity (In Pascals * Seconds)
= Strain Rate (In 1/Seconds)
Crystal-Plastic Deformation
Crystal-Plastic Deformation occurs at the atomic scale and is governed by its own set of specific mechanisms that deform crystals by the movements of atoms and atomic planes through the crystal lattice. Like viscous deformation, it is also a permanent form of deformation. Mechanisms of crystal-plastic deformation include Pressure solution, Dislocation creep, and Diffusion creep.
Biological materials
In addition to rocks, biological materials such as wood, lumber, bone, etc. can be assessed for their ductility as well, for many behave in the same manner and possess the same characteristics as abiotic Earth materials. This assessment was done in Hiroshi Yoshihara's experiment, "Plasticity Analysis of the Strain in the Tangential Direction of Solid Wood Subjected to Compression Load in the Longitudinal Direction." The study aimed to analyze the behavioral rheology of 2 wood specimens, the Sitka Spruce and Japanese Birch. In the past, it was shown that solid wood, when subjected to compressional stresses, initially has a linear stress-strain diagram (indicative of elastic deformation) and later, under greater load, demonstrates a non-linear diagram indicative of ductile objects. To analyze the rheology, the stress was restricted to uniaxial compression in the longitudinal direction and the post-linear behavior was analyzed using plasticity theory. Controls included moisture content in the lumber, lack of defects such as knots or grain distortions, temperature at 20 C, relative humidity at 65%, and size of the cut shapes of the wood samples.
Results obtained from the experiment exhibited a linear stress-strain relationship during elastic deformation but also an unexpected non-linear relationship between stress and strain for the lumber after the elastic limit was reached, deviating from the model of plasticity theory. Multiple reasons were suggested as to why this came about. First, since wood is a biological material, it was suggested that under great stress in the experiment, the crushing of cells within the sample could have been a cause for deviation from perfectly plastic behavior. With greater destruction of cellular material, the stress-strain relationship is hypothesized to become more and more nonlinear and non-ideal with greater stress. Additionally, because the samples were inhomogeneous (non-uniform) materials, it was assumed that some bending or distortion may have occurred in the samples that could have deviated the stress from being perfectly uniaxial. This may have also been induced by other factors like irregularities in the cellular density profile and distorted sample cutting.
The conclusions of the research accurately showed that although biological materials can behave like rocks undergoing deformation, there are many other factors and variables that must be considered, making it difficult to standardize the ductility and material properties of a biological substance.
Peak Ductility Demand
Peak Ductility Demand is a quantity used particularly in the fields of architecture, geological engineering, and mechanical engineering. It is defined as the amount of ductile deformation a material must be able to withstand (when exposed to a stress) without brittle fracture or failure. This quantity is particularly useful in the analysis of failure of structures in response to earthquakes and seismic waves.
It has been shown that earthquake aftershocks can increase the peak ductility demand with respect to the mainshocks by up to 10%.
| Physical sciences | Structural geology | Earth science |
6832171 | https://en.wikipedia.org/wiki/Saurischia | Saurischia | Saurischia ( , meaning "reptile-hipped" from the Greek () meaning 'lizard' and () meaning 'hip joint') is one of the two basic divisions of dinosaurs (the other being Ornithischia), classified by their hip structure. Saurischia and Ornithischia were originally called orders by Harry Seeley in 1888 though today most paleontologists classify Saurischia as an unranked clade rather than an order.
Description
All carnivorous dinosaurs (certain types of theropods) are traditionally classified as saurischians, as are all of the birds and one of the two primary lineages of herbivorous dinosaurs, the sauropodomorphs. At the end of the Cretaceous Period, all saurischians except birds became extinct in the course of the Cretaceous–Paleogene extinction event. Birds, as a group of maniraptoran theropod dinosaurs, are a sub-clade of saurischian dinosaurs in phylogenetic classification.
Saurischian dinosaurs are traditionally distinguished from ornithischian dinosaurs by their three-pronged pelvic structure, with the pubis pointed forward. The ornithischians' pelvis is arranged with the pubis rotated backward, parallel with the ischium, often also with a forward-pointing process, giving a four-pronged structure. The saurischian hip structure led Seeley to name them "lizard-hipped" dinosaurs, because they retained the ancestral hip anatomy also found in modern lizards and other reptiles. He named ornithischians "bird-hipped" dinosaurs because their hip arrangement was superficially similar to that of birds, though he did not propose any specific relationship between ornithischians and birds. However, in the view which has long been held, this "bird-hipped" arrangement evolved several times independently in dinosaurs, first in the ornithischians, then in the lineage of saurischians including birds (Avialae), and lastly in the therizinosaurians. This would then be an example of convergent evolution: avialans, therizinosaurians, and ornithischian dinosaurs all developed a similar hip anatomy independently of each other, possibly as an adaptation to their herbivorous or omnivorous diets.
Classification
In his paper naming the two groups, Seeley reviewed previous classification schemes put forth by other paleontologists to divide up the traditional order Dinosauria. He preferred one that had been put forward by Othniel Charles Marsh in 1878, which divided dinosaurs into four orders: Sauropoda, Theropoda, Ornithopoda, and Stegosauria (these names are still used today in much the same way to refer to suborders or clades within Saurischia and Ornithischia).
Seeley, however, wanted to formulate a classification that would take into account a single primary difference between major dinosaurian groups based on a characteristic that also differentiated them from other reptiles. He found this in the configuration of the hip bones, and found that all four of Marsh's orders could be divided neatly into two major groups based on this feature. He placed the Stegosauria and Ornithopoda in the Ornithischia, and the Theropoda and Sauropoda in the Saurischia. Furthermore, Seeley used this major difference in the hip bones, along with many other noted differences between the two groups, to argue that "dinosaurs" were not a natural grouping at all, but rather two distinct orders that had arisen independently from more primitive archosaurs. This concept that "dinosaur" was an outdated term for two distinct orders lasted many decades in the scientific and popular literature, and it was not until the 1960s that scientists began to again consider the possibility that saurischians and ornithischians were more closely related to each other than they were to other archosaurs.
Although his concept of a polyphyletic Dinosauria is no longer accepted by most paleontologists, Seeley's basic division of the two dinosaurian groups has stood the test of time, and has been supported by modern cladistic analysis of relationships among dinosaurs. A node-base clade, Eusaurischia, was named for the least inclusive group containing sauropodomorphs (represented by Cetiosaurus) and theropods (represented by Neornithes). Any saurischian that diverged before the theropod-sauropodomorph split is therefore outside clade Eusaurischia.
One alternative hypothesis challenging Seeley's classification was proposed by Robert T. Bakker in his 1986 book The Dinosaur Heresies. Bakker's classification separated the theropods into their own group and placed the two groups of herbivorous dinosaurs (the sauropodomorphs and ornithischians) together in a separate group he named the Phytodinosauria ("plant dinosaurs"). The Phytodinosauria hypothesis was based partly on the supposed link between ornithischians and prosauropods, and the idea that the former had evolved directly from the latter, possibly by way of an enigmatic family that seemed to possess characters of both groups, the segnosaurs. However, it was later found that segnosaurs were an unusual type of herbivorous theropod saurischian closely related to birds, and the Phytodinosauria hypothesis fell out of favor.
A 2017 study by Matthew Grant Baron, David B. Norman and Paul M. Barrett did not find support for a monophyletic Saurischia, according to its traditional definition. Instead, the group was found to be paraphyletic. As a solution, Theropoda was removed from the group and placed as the sister group to the Ornithischia in the newly defined clade Ornithoscelida. As another result, the authors redefined Saurischia as "the most inclusive clade that contains D[iplodocus] carnegii, but not T[riceratops] horridus", resulting in a clade containing only the Sauropodomorpha and Herrerasauridae. Thomas Holtz (2017) recommended using the name Sauropodomorpha to refer to a possible clade that includes traditional sauropodomorphs and herrerasaurids; alternatively, he proposed redefining the long-disused taxon Pachypodosauria to include Sauropodomorpha and Herrerasauridae as subclades. Cau (2018) also supported Ornithoscelida but placed herrerasaurids, Tawa and Daemonosaurus in a clade (Herrerasauria) outside Dinosauria. Other recent studies support a view closer to the traditional Saurischia hypothesis, with theropods closer to sauropodomorphs than to ornithischians. Novas et al. (2021) support Cau's herrerasaur phylogeny but place this clade in Saurischia.
| Biology and health sciences | Dinosaurs | Animals |
3853380 | https://en.wikipedia.org/wiki/Stem-cell%20therapy | Stem-cell therapy | Stem-cell therapy uses stem cells to treat or prevent a disease or condition. , the only FDA-approved therapy using stem cells is hematopoietic stem cell transplantation. This usually takes the form of a bone marrow or peripheral blood stem cell transplantation, but the cells can also be derived from umbilical cord blood. Research is underway to develop various sources for stem cells as well as to apply stem-cell treatments for neurodegenerative diseases and conditions such as diabetes and heart disease.
Stem-cell therapy has become controversial following developments such as the ability of scientists to isolate and culture embryonic stem cells, to create stem cells using somatic cell nuclear transfer, and their use of techniques to create induced pluripotent stem cells. This controversy is often related to abortion politics and human cloning. Additionally, efforts to market treatments based on transplant of stored umbilical cord blood have been controversial.
Medical uses
For over 90 years, hematopoietic stem cell transplantation (HSCT) has been used to treat people with conditions such as leukemia and lymphoma; this is the only widely practiced form of stem-cell therapy. During chemotherapy, most growing cells are killed by the cytotoxic agents. These agents, however, cannot discriminate between the leukaemia or neoplastic cells, and the hematopoietic stem cells within the bone marrow. This is the side effect of conventional chemotherapy strategies that the stem-cell transplant attempts to reverse; a donor's healthy bone marrow reintroduces functional stem cells to replace the cells lost in the host's body during treatment. The transplanted cells also generate an immune response that helps to kill off the cancer cells; this process can go too far, however, leading to graft vs host disease, the most serious side effect of this treatment.
Another stem-cell therapy, called Prococvhymal, was conditionally approved in Canada in 2012 for the management of acute graft-vs-host disease in children who are unresponsive to steroids. It is an allogenic stem therapy based on mesenchymal stem cells (MSCs) derived from the bone marrow of adult donors. MSCs are purified from the marrow, cultured and packaged, with up to 10,000 doses derived from a single donor. The doses are stored frozen until needed.
The FDA has approved five hematopoietic stem-cell products derived from umbilical cord blood, for the treatment of blood and immunological diseases.
In 2014, the European Medicines Agency recommended approval of limbal stem cells for people with severe limbal stem cell deficiency due to burns in the eye.
Research only
Stem cells are being studied for several reasons. The molecules and exosomes released from stem cells are also being studied in an effort to make medications. In addition to the functions of the cells themselves, paracrine soluble factors produced by stem cells, known as the stem cell secretome, have been found to be another mechanism by which stem cell-based therapies mediate their effects in degenerative, autoimmune, and inflammatory diseases.
Sources for human stem cells
Most stem cells intended for regenerative therapy are generally isolated either from the patient's bone marrow or from adipose tissue. Mesenchymal stem cells can differentiate into the cells that make up bone, cartilage, tendons, and ligaments, as well as muscle, neural and other progenitor tissues. They have been the main type of stem cells studied in the treatment of diseases affecting these tissues. The number of stem cells transplanted into damaged tissue may alter the efficacy of treatment. Accordingly, stem cells derived from bone marrow aspirates, for instance, are cultured in specialized laboratories for expansion to millions of cells. Although adipose-derived tissue also requires processing prior to use, the culturing methodology for adipose-derived stem cells is not as extensive as that for bone marrow-derived cells. While it is thought that bone-marrow-derived stem cells are preferred for bone, cartilage, ligament, and tendon repair, others believe that the less challenging collection techniques and the multi-cellular microenvironment already present in adipose-derived stem cell fractions make the latter the preferred source for autologous transplantation.
New sources of mesenchymal stem cells are being researched, including stem cells present in the skin and dermis which are of interest because of the ease at which they can be harvested with minimal risk to the animal. Hematopoietic stem cells have also been discovered to be travelling in the blood stream and possess equal differentiating ability as other mesenchymal stem cells, again with a very non-invasive harvesting technique.
There has been more recent interest in the use of extra embryonic mesenchymal stem cells. Research is underway to examine the differentiating capabilities of stem cells found in the umbilical cord, yolk sac and placenta of different animals. These stem cells are thought to have more differentiating ability than their adult counterparts, including the ability to more readily form tissues of endodermal and ectodermal origin.
Embryonic stem cell lines
As of 2010, there was widespread controversy over the use of human embryonic stem cells. This controversy primarily targets the techniques used to derive new embryonic stem cell lines, which often requires the destruction of the blastocyst. Opposition to the use of human embryonic stem cells in research is often based on philosophical, moral, or religious objections. There is other stem cell research that does not involve the destruction of a human embryo, and such research involves adult stem cells, amniotic stem cells, and induced pluripotent stem cells.
In January 2009, the US Food and Drug Administration gave clearance to Geron Corporation for the first clinical trial of an embryonic stem-cell-based therapy on humans. The trial aimed to evaluate the drug GRNOPC1, embryonic stem cell-derived oligodendrocyte progenitor cells, on people with acute spinal cord injury. The trial was discontinued in November 2011 so that the company could focus on therapies in the "current environment of capital scarcity and uncertain economic conditions". In 2013 biotechnology and regenerative medicine company BioTime () acquired Geron's stem cell assets in a stock transaction, with the aim of restarting the clinical trial.
Mesenchymal stromal cells (MSCs)
Scientists reported 2012 that MSCs when transfused immediately within few hours post thawing may show reduced function or show decreased efficacy in treating diseases as compared to those MSCs which are in log phase of cell growth (fresh), so cryopreserved MSCs should be brought back into log phase of cell growth in in vitro culture before administration. Re-culturing of MSCs will help in recovering from the shock the cells get during freezing and thawing. Various MSC clinical trials which used cryopreserved product immediately post thaw have failed as compared to those clinical trials which used fresh MSCs.
In drug discovery and biomedical research
The ability to grow up functional adult tissues indefinitely in culture through Directed differentiation creates new opportunities for drug research. Researchers are able to grow up differentiated cell lines and then test new drugs on each cell type to examine possible interactions in vitro before performing in vivo studies. This is critical in the development of drugs for use in veterinary research because of the possibilities of species-specific interactions. The hope is that having these cell lines available for research use will reduce the need for research animals used because effects on human tissue in vitro will provide insight not normally known before the animal testing phase.
With the advent of induced pluripotent stem cells (iPSC), treatments being explored and created for the used in endangered low production animals possible. Rather than needing to harvest embryos or eggs, which are limited, the researchers can remove mesenchymal stem cells with greater ease and greatly reducing the danger to the animal due to noninvasive techniques. This allows the limited eggs to be put to use for reproductive purposes only.
Stem cell expansion
To be used for research or treatment applications, large numbers of high-quality stem cells are needed. Thus, it is necessary to develop culture systems which produce pure populations of tissue-specific stem cells in vitro without the loss of stem-cell potential. Two main approaches are taken for this purpose: two-dimensional and three-dimensional cell culture.
Cell culture in two dimensions has been routinely performed in thousands of laboratories worldwide for the past four decades. In two-dimensional platforms, cells are typically exposed to a solid, rigid flat surface on the basal side and to liquid at the apical
surface. Inhabiting such a two-dimensional rigid substrate requires a dramatic adaption for the surviving cells because they lack the extracellular matrix that is unique to each cell type which may alter cell metabolism and reduce its functionality.
Three-dimensional cell culture systems may create a biomimicking microenvironment for stem cells, resembling their native three-dimensional extracellular matrix (ECM). Advanced biomaterials have significantly contributed to three-dimensional cell culture systems in recent decades, and more unique and complex biomaterials have been proposed for improving stem-cell proliferation and controlled differentiation. Among them, nanostructured biomaterials are of particular interest because they have the advantage of a high surface-to-volume ratio, and they mimic the physical and biological features of natural ECM at the nanoscale.
Assumed regenerative models
Stem cells are thought to mediate repair via five primary mechanisms: 1) providing an anti-inflammatory effect, 2) homing to damaged tissues and recruiting other cells, such as endothelial progenitor cells, that are necessary for tissue growth, 3) supporting tissue remodeling over scar formation, 4) inhibiting apoptosis, and 5) differentiating into bone, cartilage, tendon, and ligament tissue.
To further enrich blood supply to the damaged areas, and consequently promote tissue regeneration, platelet-rich plasma could be used in conjunction with stem cell transplantation. The efficacy of some stem cell populations may also be affected by the method of delivery; for instance, to regenerate bone, stem cells are often introduced in a scaffold where they produce the minerals necessary for generation of functional bone.
Stem cells have been shown to have low immunogenicity due to the relatively low number of MHC molecules found on their surface. In addition, they have been found to secrete chemokines that alter the immune response and promote tolerance of the new tissue. This allows for allogeneic treatments to be performed without a high rejection risk.
Potential applications
Neurodegeneration
Research has been conducted on the effects of stem cells on animal models of brain degeneration, such as in Parkinson's disease, Amyotrophic lateral sclerosis, and Alzheimer's disease. Preliminary studies related to multiple sclerosis have been conducted, and a 2020 phase 2 trial found significantly improved outcomes for mesenchymal stem cell treated patients compared to those receiving a sham treatment. In January 2021 the FDA approved the first clinical trial for an investigational stem cell therapy to restore lost brain cells in people with advanced Parkinson's disease.
Healthy adult brains contain neural stem cells, which divide to maintain general stem-cell numbers, or become progenitor cells. In healthy adult laboratory animals, progenitor cells migrate within the brain and function primarily to maintain neuron populations for olfaction (the sense of smell). Pharmacological activation of endogenous neural stem cells has been reported to induce neuroprotection and behavioral recovery in adult rat models of neurological disorders.
Brain and spinal cord injury
Stroke and traumatic brain injury lead to cell death, characterized by a loss of neurons and oligodendrocytes within the brain. Clinical and animal studies have been conducted into the experimental use of stem cells in cases of spinal cord injury.
Frailty syndrome
In 2017, a small-scale study on individuals 60 years or older with aging frailty showed, after intravenous treatment with Mesenchymal stem cells (MSC) from healthy young donors, significant improvements in physical performance measures. MSC helps with the blockade of inflammation by decreasing it, causing the effects of frailty to reverse.
Heart
In 2012, stem cells were studied in people with severe heart disease. The work by Bodo-Eckehard Strauer was discredited by identifying hundreds of factual contradictions. Among several clinical trials reporting that adult stem cell therapy is safe and effective, actual evidence of benefit has been reported from only a few studies. Some preliminary clinical trials achieved only modest improvements in heart function following the use of bone marrow stem cell therapy.
Stem-cell therapy for the treatment of myocardial infarction usually makes use of autologous bone marrow stem cells, but other types of adult stem cells may be used, such as adipose-derived stem cells.
Possible mechanisms of recovery include:Generation of heart muscle cells, Stimulating growth of new blood vessels to repopulate damaged heart tissue and secretion of growth factors.
Blood-cell formation
The specificity of the human immune-cell repertoire is what allows the human body to defend itself from rapidly adapting antigens. However, the immune system is vulnerable to degradation upon the pathogenesis of disease, and because of the critical role that it plays in overall defense, its degradation is often fatal to the organism as a whole. Diseases of hematopoietic cells are diagnosed and classified via a subspecialty of pathology known as hematopathology. The specificity of the immune cells is what allows recognition of foreign antigens, causing further challenges in the treatment of immune disease. Identical matches between donor and recipient are no longer required for successful transplantation. Rather, haploidentical matches have facilitated numerous transplants, given improvements in post-transplant immunosuppressive regimens. Research using both hematopoietic adult stem cells and embryonic stem cells has provided insight into the possible mechanisms and methods of treatment for many of these ailments.
Fully mature human red blood cells may be generated ex vivo by hematopoietic stem cells (HSCs), which are precursors of red blood cells. In this process, HSCs are grown together with stromal cells, creating an environment that mimics the conditions of bone marrow, the natural site of red-blood-cell growth. Erythropoietin, a growth factor, is added, coaxing the stem cells to complete terminal differentiation into red blood cells. Further research into this technique should have potential benefits for gene therapy, blood transfusion, and topical medicine.
Regrowing teeth
In 2004, scientists at King's College London discovered a way to cultivate a complete tooth in mice and were able to grow bioengineered teeth stand-alone in the laboratory. Researchers are confident that tooth regeneration technology can be used to grow live teeth in people.
In theory, stem cells taken from the patient could be coaxed in the lab turning into a tooth bud which, when implanted in the gums, will give rise to a new tooth, and would be expected to be grown in a time over three weeks. It will fuse with the jawbone and release chemicals that encourage nerves and blood vessels to connect with it. The process is similar to what happens when humans grow their original adult teeth. Many challenges remain, however, before stem cells can be a choice for the replacement of missing teeth in the future.
Cochlear hair cell regrowth
Heller has reported success in re-growing cochlea hair cells with the use of embryonic stem cells.
In a 2019 review that looked at hearing regeneration and regenerative medicine, stem cell-derived otic progenitors have the potential to greatly improve hearing.
Blindness and vision impairment
Since 2003, researchers have successfully transplanted corneal stem cells into damaged eyes to restore vision. "Sheets of retinal cells used by the team are harvested from aborted fetuses, which some people find objectionable." When these sheets are transplanted over the damaged cornea, the stem cells stimulate renewed repair, eventually restoring vision. The latest such development was in June 2005, when researchers at the Queen Victoria Hospital of Sussex, England were able to restore the sight of forty people using the same technique. The group, led by Sheraz Daya, was able to successfully use adult stem cells obtained from the patient, a relative, or even a cadaver. Further rounds of trials are ongoing.
Pancreatic beta cells
People with Type 1 diabetes lose the function of insulin-producing beta cells within the pancreas. In a 2007 publication of experiments, scientists have been able to coax embryonic stem cells to turn into beta cells in the lab. In theory, if the beta cell is transplanted successfully, they will be able to replace malfunctioning ones in a diabetic patient. There are adverse effects of high glucose concentrations on stem cell therapy, however.
Orthopedics
As of 2017, use of mesenchymal stem cells (MSCs) derived from adult stem cells was under preliminary research for potential orthopedic applications in bone and muscle trauma, cartilage repair, osteoarthritis, intervertebral disc surgery, rotator cuff surgery, and musculoskeletal disorders, among others. Other areas of orthopedic research for uses of MSCs include tissue engineering and regenerative medicine.
Wound healing
Stem cells can also be used to stimulate the growth of human tissues. In an adult, wounded tissue is most often replaced by scar tissue, which is characterized in the skin by disorganized collagen structure, loss of hair follicles and irregular vascular structure. In the case of wounded fetal tissue, however, wounded tissue is replaced with normal tissue through the activity of stem cells. A possible method for tissue regeneration in adults is to place adult stem cell "seeds" inside a tissue bed "soil" in a wound bed and allow the stem cells to stimulate differentiation in the tissue bed cells. This method elicits a regenerative response more similar to fetal wound-healing than adult scar tissue formation. As of 2018, researchers were still investigating different aspects of the "soil" tissue that are conducive to regeneration. Because of the general healing capabilities of stem cells, they have gained interest for the treatment of cutaneous wounds, such as in skin cancer.
HIV/AIDS
In 2013, scientists have been investigating an alternative approach to treating HIV-1/AIDS, based on the creation of a disease-resistant immune system through transplantation of autologous, gene-modified (HIV-1-resistant) hematopoietic stem and progenitor cells (GM-HSPC).
Criticisms
In 2013, studies of autologous bone marrow stem cells on ventricular function were found to contain "hundreds" of discrepancies. Critics report that of 48 reports, just five underlying trials seemed to be used, and that in many cases whether they were randomized or merely observational accepter-versus-rejecter, was contradictory between reports of the same trial. One pair of reports of identical baseline characteristics and final results, was presented in two publications as, respectively, a 578-patient randomized trial and as a 391-subject observational study. Other reports required (impossible) negative standard deviations in subsets of people or contained fractional subjects, negative NYHA classes. Overall, many more people were reported as having receiving stem cells in trials, than the number of stem cells processed in the hospital's laboratory during that time. A university investigation, closed in 2012 without reporting, was reopened in July 2013.
In 2014, a meta-analysis on stem cell therapy using bone-marrow stem cells for heart disease revealed discrepancies in published clinical trial reports, whereby studies with a higher number of discrepancies showed an increase in effect sizes. Another meta-analysis based on the intra-subject data of 12 randomized trials was unable to find any significant benefits of stem cell therapy on primary endpoints, such as major adverse events or increase in heart function measures, concluding there was no benefit.
2018 results of the TIME trial, which used a randomized, double-blind, placebo-controlled trial design, concluded that "bone marrow mononuclear cells administration did not improve recovery of LV function over 2 years" in people who had a myocardial infarction. Accordingly, the BOOST-2 trial conducted in 10 medical centers in Germany and Norway reported that the trial result "does not support the use of nucleated BMCs in patients with STEMI and moderately reduced LVEF". Furthermore, the trial also did not meet any other secondary MRI endpoints, leading to a conclusion that intracoronary bone marrow stem cell therapy does not offer a functional or clinical benefit.
In 2021, stem cell injections in the US have caused grave infections in at least 20 patients who received umbilical cord blood-derived products marketed as "stem cell treatment". In 2023, the case of a woman who was infected with Mycobacterium abscessus and sustained meningitis after stem cell treatment for multiple sclerosis at a commercial clinic in Baja California, Mexico was published.
Veterinary medicine
Research conducted on horses, dogs, and cats has led to the development of stem cell treatments in veterinary medicine which can target a wide range of injuries and diseases, such as myocardial infarction, stroke, tendon and ligament damage, osteoarthritis, osteochondrosis and muscular dystrophy, both in large animals as well as in humans. While investigation of cell-based therapeutics generally reflects human medical needs, the high degree of frequency and severity of certain injuries in racehorses has put veterinary medicine at the forefront of this novel regenerative approach. Companion animals can serve as clinically relevant models that closely mimic human disease.
Sources of veterinarian stem cells
Veterinary applications of stem cell therapy as a means of tissue regeneration have been largely shaped by research that began with the use of adult-derived mesenchymal stem cells to treat animals with injuries or defects affecting bone, cartilage, ligaments and/or tendons. There are two main categories of stem cells used for treatments: allogeneic stem cells derived from a genetically different donor within the same species, and autologous mesenchymal stem cells, derived from the patient before use in various treatments. A third category, xenogenic stem cells, or stem cells derived from different species, are used primarily for research purposes, especially for human treatments.
Bone repair
Bone has a unique and well-documented natural healing process that normally is sufficient to repair fractures and other common injuries. Misaligned breaks due to severe trauma, as well as treatments like tumor resections of bone cancer, are prone to improper healing if left to the natural process alone. Scaffolds composed of natural and artificial components are seeded with mesenchymal stem cells and placed in the defect. Within four weeks of placing the scaffold, newly formed bone begins to integrate with the old bone and within 32 weeks, full union is achieved. Further studies are necessary to fully characterize the use of cell-based therapeutics for treatment of bone fractures.
Stem cells have been used to treat degenerative bone diseases in dogs. The normally recommended treatment for dogs that have Legg–Calve–Perthes disease is to remove the head of the femur after the degeneration has progressed. Recently, mesenchymal stem cells have been injected directly in to the head of the femur, with success not only in bone regeneration, but also in pain reduction.
Ligament and tendon repair
Autologous stem cell-based treatments for ligament injury, tendon injury, osteoarthritis, osteochondrosis, and sub-chondral bone cysts have been commercially available to practicing veterinarians to treat horses since 2003 in the United States and since 2006 in the United Kingdom. Autologous stem cell based treatments for tendon injury, ligament injury, and osteoarthritis in dogs have been available to veterinarians in the United States since 2005. Over 3000 privately owned horses and dogs have been treated with autologous adipose-derived stem cells. The efficacy of these treatments has been shown in double-blind clinical trials for dogs with osteoarthritis of the hip and elbow and horses with tendon damage.
Race horses are especially prone to injuries of the tendon and ligaments. Conventional therapies are very unsuccessful in returning the horse to full functioning potential. Natural healing, guided by the conventional treatments, leads to the formation of fibrous scar tissue that reduces flexibility and full joint movement. Traditional treatments prevented a large number of horses from returning to full activity and also have a high incidence of re-injury due to the stiff nature of the scarred tendon. Introduction of both bone marrow and adipose derived stem cells, along with natural mechanical stimulus promoted the regeneration of tendon tissue. The natural movement promoted the alignment of the new fibers and tendocytes with the natural alignment found in uninjured tendons. Stem cell treatment not only allowed more horses to return to full duty and also greatly reduced the re-injury rate over a three-year period.
The use of embryonic stem cells has also been applied to tendon repair. The embryonic stem cells were shown to have a better survival rate in the tendon as well as better migrating capabilities to reach all areas of damaged tendon. The overall repair quality was also higher, with better tendon architecture and collagen formed. There was also no tumor formation seen during the three-month experimental period. Long-term studies need to be carried out to examine the long-term efficacy and risks associated with the use of embryonic stem cells. Similar results have been found in small animals.
Joint repair
Osteoarthritis is the main cause of joint pain both in animals and humans. Horses and dogs are most frequently affected by arthritis. Natural cartilage regeneration is very limited. Different types of mesenchymal stem cells and other additives are still being researched to find the best type of cell and method for long-term treatment.
Adipose-derived mesenchymal cells are currently the most often used for stem cell treatment of osteoarthritis because of the non-invasive harvesting. This is a recently developed, non-invasive technique developed for easier clinical use. Dogs receiving this treatment showed greater flexibility in their joints and less pain.
Muscle repair
Stem cells have successfully been used to ameliorate healing in the heart after myocardial infarction in dogs. Adipose and bone marrow derived stem cells were removed and induced to a cardiac cell fate before being injected into the heart. The heart was found to have improved contractility and a reduction in the damaged area four weeks after the stem cells were applied.
In 2007, a trial was underway for a patch made of a porous substance onto which the stem cells are "seeded" in order to induce tissue regeneration in heart defects. Tissue was regenerated and the patch was well incorporated into the heart tissue. This is thought to be due, in part, to improved angiogenesis and reduction of inflammation. Although cardiomyocytes were produced from the mesenchymal stem cells, they did not appear to be contractile. Other treatments that induced a cardiac fate in the cells before transplanting had greater success at creating contractile heart tissue.
2018 research, such as the European nTRACK research project, aims to demonstrate that multimodal nanoparticles can structurally and functionally track stem cell in muscle regeneration therapy. The idea is to label stem cells with gold nano-particles that are fully characterised for uptake, functionality, and safety. The labelled stem cells will be injected into an injured muscle and tracked using imaging systems. However, the system still needs to be demonstrated at lab scale.
Nervous system repair
Spinal cord injuries are one of the most common traumas brought into veterinary hospitals. Spinal injuries occur in two ways after the trauma: the primary mechanical damage, and in secondary processes, like inflammation and scar formation, in the days following the trauma. These cells involved in the secondary damage response secrete factors that promote scar formation and inhibit cellular regeneration. Mesenchymal stem cells that are induced to a neural cell fate are loaded onto a porous scaffold and are then implanted at the site of injury. The cells and scaffold secrete factors that counteract those secreted by scar forming cells and promote neural regeneration. Eight weeks later, dogs treated with stem cells showed immense improvement over those treated with conventional therapies. Dogs treated with stem cells were able to occasionally support their own weight, which has not been seen in dogs undergoing conventional therapies.
In a study to evaluate the treatment of experimentally induced MS in dogs using laser activated non-expanded adipose derived stem cells. The results showed amelioration of the clinical signs over time confirmed by the resolution of the previous lesions on MRI. Positive migration of the injected cells to the site of lesion, increased remyelination detected by Myelin Basic Proteins, positive differentiation into Olig2 positive oligodendrocytes, prevented the glial scar formation and restored axonal architecture.
Treatments are also in clinical trials to repair and regenerate peripheral nerves. Peripheral nerves are more likely to be damaged, but the effects of the damage are not as widespread as seen in injuries to the spinal cord. Treatments are currently in clinical trials to repair severed nerves, with early success. Stem cells induced to a neural fate injected in to a severed nerve. Within four weeks, regeneration of previously damaged stem cells and completely formed nerve bundles were observed.
Stem cells are also in clinical phases for treatment in ophthalmology. Hematopoietic stem cells have been used to treat corneal ulcers of different origin of several horses. These ulcers were resistant to conventional treatments available, but quickly responded positively to the stem cell treatment. Stem cells were also able to restore sight in one eye of a horse with retinal detachment, allowing the horse to return to daily activities.
Conservation
Stem cells are being explored for use in conservation efforts. Spermatogonial stem cells have been harvested from a rat and placed into a mouse host and fully mature sperm were produced with the ability to produce viable offspring. Currently research is underway to find suitable hosts for the introduction of donor spermatogonial stem cells. If this becomes a viable option for conservationists, sperm can be produced from high genetic quality individuals who die before reaching sexual maturity, preserving a line that would otherwise be lost.
Society and culture
Marketing and costs
In the late 1990s and early 2000s, there was an initial wave of companies and clinics offering stem cell therapy, while not substantiating health claims or having regulatory approval. By 2012, a second wave of companies and clinics had emerged, usually located in developing countries where medicine is less regulated and offering stem cell therapies on a medical tourism model. Like the first wave companies and clinics, they made similar strong, but unsubstantiated, claims, mainly by clinics in the United States, Mexico, Thailand, India, and South Africa. By 2016, research indicated that there were more than 550 stem cell clinics in the US alone selling generally unproven therapies for a wide array of medical conditions in almost every state in the country, altering the dynamic of stem cell tourism. In 2018, the FDA sent a warning letter to StemGenex Biologic Laboratories in San Diego, which marketed a service in which it took body fat from people, processed it into mixtures it said contained various forms of stem cells, and administered it back to the person by inhalation, intravenously, or infusion into their spinal cords; the company said the treatment was useful for many chronic and life-threatening conditions.
One common marketing tactic is registering on ClinicalTrials.gov, the US government database for clinical trials. Registration of a study notifies the agency but does not prove that review has taken place. Registration with the FDA similarly does not prove that approval has been granted.
Costs of stem cell therapies range widely by clinic, condition, and cell type, but most commonly range between $10,000-$20,000. Insurance does not cover stem cell injections at clinics so patients often use on-line fundraising. In 2018, the US Federal Trade Commission found health centers and an individual physician making unsubstantiated claims for stem cell therapies, and forced refunds of some $500,000. The FDA filed suit against two stem cell clinic firms around the same time, seeking permanent injunctions against their marketing and use of unapproved adipose stem cell products.
COVID-19-related marketing and government agency responses
Although according to the NIH no stem cell treatments have been approved for COVID-19, and the agency recommends against the use of MSCs for the disease, some stem cell clinics began marketing both unproven and non-FDA-approved stem cells and exosomes for COVID-19 in 2020. The FDA took prompt action by sending letters to the firms in question. The FTC also warned a stem cell firm for misleading COVID-19-related marketing.
| Technology | Biotechnology | null |
3854225 | https://en.wikipedia.org/wiki/Waveguide%20%28radio%20frequency%29 | Waveguide (radio frequency) | In radio-frequency engineering and communications engineering, a waveguide is a hollow metal pipe used to carry radio waves. This type of waveguide is used as a transmission line mostly at microwave frequencies, for such purposes as connecting microwave transmitters and receivers to their antennas, in equipment such as microwave ovens, radar sets, satellite communications, and microwave radio links.
The electromagnetic waves in a (metal-pipe) waveguide may be imagined as travelling down the guide in a zig-zag path, being repeatedly reflected between opposite walls of the guide. For the particular case of rectangular waveguide, it is possible to base an exact analysis on this view. Propagation in a dielectric waveguide may be viewed in the same way, with the waves confined to the dielectric by total internal reflection at its surface. Some structures, such as non-radiative dielectric waveguides and the Goubau line, use both metal walls and dielectric surfaces to confine the wave.
Principle
Depending on the frequency, waveguides can be constructed from either conductive or dielectric materials. Generally, the lower the frequency to be passed the larger the waveguide is. For example, the natural waveguide the earth forms given by the dimensions between the conductive ionosphere and the ground as well as the circumference at the median altitude of the Earth is resonant at 7.83 Hz. This is known as Schumann resonance. On the other hand, waveguides used in extremely high frequency (EHF) communications can be less than a millimeter in width.
History
During the 1890s theorists did the first analyses of electromagnetic waves in ducts. Around 1893 J. J. Thomson derived the electromagnetic modes inside a cylindrical metal cavity. In 1897 Lord Rayleigh did a definitive analysis of waveguides; he solved the boundary value problem of electromagnetic waves propagating through both conducting tubes and dielectric rods of arbitrary shape. He showed that the waves could travel without attenuation only in specific normal modes with either the electric field (TE modes) or magnetic field (TM modes), perpendicular to the direction of propagation. He also showed each mode had a cutoff frequency below which waves would not propagate. Since the cutoff wavelength for a given tube was of the same order as its width, it was clear that a hollow conducting tube could not carry radio wavelengths much larger than its diameter. In 1902 R. H. Weber observed that electromagnetic waves travel at a slower speed in tubes than in free space, and deduced the reason; that the waves travel in a "zigzag" path as they reflect from the walls.
Prior to the 1920s, practical work on radio waves concentrated on the low frequency end of the radio spectrum, as these frequencies were better for long-range communication. These were far below the frequencies that could propagate in even large waveguides, so there was little experimental work on waveguides during this period, although a few experiments were done. In a June 1, 1894 lecture, "The work of Hertz", before the Royal Society, Oliver Lodge demonstrated the transmission of 3 inch radio waves from a spark gap through a short cylindrical copper duct. In his pioneering 1894-1900 research on microwaves, Jagadish Chandra Bose used short lengths of pipe to conduct the waves, so some sources credit him with inventing the waveguide. However, after this, the concept of radio waves being carried by a tube or duct passed out of engineering knowledge.
During the 1920s the first continuous sources of high frequency radio waves were developed: the Barkhausen–Kurz tube, the first oscillator which could produce power at UHF frequencies; and the split-anode magnetron which by the 1930s had generated radio waves at up to 10 GHz. These made possible the first systematic research on microwaves in the 1930s. It was discovered that transmission lines used to carry lower frequency radio waves, parallel line and coaxial cable, had excessive power losses at microwave frequencies, creating a need for a new transmission method.
The waveguide was developed independently between 1932 and 1936 by George C. Southworth at Bell Telephone Laboratories and Wilmer L. Barrow at the Massachusetts Institute of Technology, who worked without knowledge of one another. Southworth's interest was sparked during his 1920s doctoral work in which he measured the dielectric constant of water with a radio frequency Lecher line in a long tank of water. He found that if he removed the Lecher line, the tank of water still showed resonance peaks, indicating it was acting as a dielectric waveguide. At Bell Labs in 1931 he resumed work in dielectric waveguides. By March 1932 he observed waves in water-filled copper pipes. Rayleigh's previous work had been forgotten, and Sergei A. Schelkunoff, a Bell Labs mathematician, did theoretical analyses of waveguides and rediscovered waveguide modes. In December 1933 it was realized that with a metal sheath the dielectric is superfluous and attention shifted to metal waveguides.
Barrow had become interested in high frequencies in 1930 studying under Arnold Sommerfeld in Germany. At MIT beginning in 1932 he worked on high frequency antennas to generate narrow beams of radio waves to locate aircraft in fog. He invented a horn antenna and hit on the idea of using a hollow pipe as a feedline to feed radio waves to the antenna. By March 1936 he had derived the propagation modes and cutoff frequency in a rectangular waveguide. The source he was using had a large wavelength of 40 cm, so for his first successful waveguide experiments he used a 16-foot section of air duct, 18 inches in diameter.
Barrow and Southworth became aware of each other's work a few weeks before both were scheduled to present papers on waveguides to a combined meeting of the American Physical Society and the Institute of Radio Engineers in May 1936. They amicably worked out credit sharing and patent division arrangements.
The development of centimeter radar during World War 2 and the first high power microwave tubes, the klystron (1938) and cavity magnetron (1940), resulted in the first widespread use of waveguide. Standard waveguide "plumbing" components were manufactured, with flanges on the end which could be bolted together. After the war in the 1950s and 60s waveguides became common in commercial microwave systems, such as airport radar and microwave relay networks which were built to transmit telephone calls and television programs between cities.
Description
In the microwave region of the electromagnetic spectrum, a waveguide normally consists of a hollow metallic conductor. These waveguides can take the form of single conductors with or without a dielectric coating, e.g. the Goubau line and helical waveguides. Hollow waveguides must be one-half wavelength or more in diameter in order to support one or more transverse wave modes.
Waveguides may be filled with pressurized gas to inhibit arcing and prevent multipaction, allowing higher power transmission. Conversely, waveguides may be required to be evacuated as part of evacuated systems (e.g. electron beam systems).
A slotted waveguide is generally used for radar and other similar applications. The waveguide serves as a feed path, and each slot is a separate radiator, thus forming an antenna. This structure has the capability of generating a radiation pattern to launch an electromagnetic wave in a specific relatively narrow and controllable direction.
A closed waveguide is an electromagnetic waveguide (a) that is tubular, usually with a circular or rectangular cross section, (b) that has electrically conducting walls, (c) that may be hollow or filled with a dielectric material, (d) that can support a large number of discrete propagating modes, though only a few may be practical, (e) in which each discrete mode defines the propagation constant for that mode, (f) in which the field at any point is describable in terms of the supported modes, (g) in which there is no radiation field, and (h) in which discontinuities and bends may cause mode conversion but not radiation.
The dimensions of a hollow metallic waveguide determine which wavelengths it can support, and in which modes. Typically the waveguide is operated so that only a single mode is present. The lowest order mode possible is generally selected. Frequencies below the guide's cutoff frequency will not propagate. It is possible to operate waveguides at higher order modes, or with multiple modes present, but this is usually impractical.
Waveguides are almost exclusively made of metal and mostly rigid structures. There are certain types of "corrugated" waveguides that have the ability to flex and bend but only used where essential since they degrade propagation properties. Due to propagation of energy in mostly air or space within the waveguide, it is one of the lowest loss transmission line types and highly preferred for high frequency applications where most other types of transmission structures introduce large losses. Due to the skin effect at high frequencies, electric current along the walls penetrates typically only a few micrometers into the metal of the inner surface. Since this is where most of the resistive loss occurs, it is important that the conductivity of interior surface be kept as high as possible. For this reason, most waveguide interior surfaces are plated with copper, silver, or gold.
Voltage standing wave ratio (VSWR) measurements may be taken to ensure that a waveguide is contiguous and has no leaks or sharp bends. If such bends or holes in the waveguide surface are present, this may diminish the performance of both transmitter and receiver equipment connected at either end. Poor transmission through the waveguide may also occur as a result of moisture build up which corrodes and degrades conductivity of the inner surfaces, which is crucial for low loss propagation. For this reason, waveguides are nominally fitted with microwave windows at the outer end that will not interfere with propagation but keep the elements out. Moisture can also cause fungus build up or arcing in high power systems such as radio or radar transmitters. Moisture in waveguides can typically be prevented with silica gel, a desiccant, or slight pressurization of the waveguide cavities with dry nitrogen or argon. Desiccant silica gel canisters may be attached with screw-on nibs and higher power systems will have pressurized tanks for maintaining pressure including leakage monitors. Arcing may also occur if there is a hole, tear or bump in the conducting walls, if transmitting at high power (usually 200 watts or more). Waveguide plumbing is crucial for proper waveguide performance. Voltage standing waves occur when impedance mismatches in the waveguide cause energy to reflect back in the opposite direction of propagation. In addition to limiting the effective transfer of energy, these reflections can cause higher voltages in the waveguide and damage equipment.
In practice
In practice, waveguides act as the equivalent of cables for super high frequency (SHF) systems. For such applications, it is desired to operate waveguides with only one mode propagating through the waveguide. With rectangular waveguides, it is possible to design the waveguide such that the frequency band over which only one mode propagates is as high as 2:1 (i.e. the ratio of the upper band edge to lower band edge is two). The relation between the waveguide dimensions and the lowest frequency is simple: if is the greater of its two dimensions, then the longest wavelength that will propagate is and the lowest frequency is thus
With circular waveguides, the highest possible bandwidth allowing only a single mode to propagate is only 1.3601:1.
Because rectangular waveguides have a much larger bandwidth over which only a single mode can propagate, standards exist for rectangular waveguides, but not for circular waveguides. In general (but not always), standard waveguides are designed such that
one band starts where another band ends, with another band that overlaps the two bands
the lower edge of the band is approximately 30% higher than the waveguide's cutoff frequency
the upper edge of the band is approximately 5% lower than the cutoff frequency of the next higher order mode
the waveguide height is half the waveguide width
The first condition is to allow for applications near band edges. The second condition limits dispersion, a phenomenon in which the velocity of propagation is a function of frequency. It also limits the loss per unit length. The third condition is to avoid evanescent-wave coupling via higher order modes. The fourth condition is that which allows a 2:1 operation bandwidth. Although it is possible to have a 2:1 operating bandwidth when the height is less than half the width, having the height exactly half the width maximizes the power that can propagate inside the waveguide before dielectric breakdown occurs.
Below is a table of standard waveguides. The waveguide name WR stands for waveguide rectangular, and the number is the inner dimension width of the waveguide in hundredths of an inch (0.01 inch = 0.254 mm) rounded to the nearest hundredth of an inch.
* Radio Components Standardization Committee
† For historical reasons the outside rather than the inside dimensions of these waveguides are 2:1 (with wall thickness WG6–WG10: 0.08" (2.0 mm), WG11A–WG15: 0.064" (1.6 mm), WG16–WG17: 0.05" (1.3 mm), WG18–WG28: 0.04" (1.0 mm))
For the frequencies in the table above, the main advantage of waveguides over coaxial cables is that waveguides support propagation with lower loss. For lower frequencies, the waveguide dimensions become impractically large, and for higher frequencies the dimensions become impractically small (the manufacturing tolerance becomes a significant portion of the waveguide size).
Mathematical analysis
Electromagnetic waveguides are analyzed by solving Maxwell's equations, or their reduced form, the electromagnetic wave equation, with boundary conditions determined by the properties of the materials and their interfaces. These equations have multiple solutions, or modes, which are eigenfunctions of the equation system. Each mode is characterized by a cutoff frequency below which the mode cannot exist in the guide. Waveguide propagation modes depend on the operating wavelength and polarization and the shape and size of the guide. The longitudinal mode of a waveguide is a particular standing wave pattern formed by waves confined in the cavity. The transverse modes are classified into different types:
TE modes (transverse electric) have no electric field in the direction of propagation.
TM modes (transverse magnetic) have no magnetic field in the direction of propagation.
TEM modes (transverse electromagnetic) have no electric nor magnetic field in the direction of propagation.
Hybrid modes have both electric and magnetic field components in the direction of propagation.
Waveguides with certain symmetries may be solved using the method of separation of variables. Rectangular wave guides may be solved in rectangular coordinates. Round waveguides may be solved in cylindrical coordinates.
In hollow, single conductor waveguides, TEM waves are not possible. This contrasts with two-conductor transmission lines used at lower frequencies; coaxial cable, parallel wire line and stripline, in which TEM mode is possible. Additionally, the propagating modes (i.e. TE and TM) inside the waveguide can be mathematically expressed as the superposition of two TEM waves.
The mode with the lowest cutoff frequency is termed the dominant mode of the guide. It is common to choose the size of the guide such that only this one mode can exist in the frequency band of operation. In rectangular and circular (hollow pipe) waveguides, the dominant modes are designated the TE1,0 mode and TE1,1 modes respectively.
Dielectric waveguides
A dielectric waveguide employs a solid dielectric rod rather than a hollow pipe. An optical fibre is a dielectric guide designed to work at optical frequencies. Transmission lines such as microstrip, coplanar waveguide, stripline or coaxial cable may also be considered to be waveguides.
Dielectric rod and slab waveguides are used to conduct radio waves, mostly at millimeter wave frequencies and above. These confine the radio waves by total internal reflection from the step in refractive index due to the change in dielectric constant at the material surface. At millimeter wave frequencies and above, metal is not a good conductor, so metal waveguides can have increasing attenuation. At these wavelengths dielectric waveguides can have lower losses than metal waveguides. Optical fibre is a form of dielectric waveguide used at optical wavelengths.
One difference between dielectric and metal waveguides is that at a metal surface the electromagnetic waves are tightly confined; at high frequencies the electric and magnetic fields penetrate a very short distance into the metal. In contrast, the surface of the dielectric waveguide is an interface between two dielectrics, so the fields of the wave penetrate outside the dielectric in the form of an evanescent (non-propagating) wave.
| Physical sciences | Electromagnetic radiation | Physics |
21464996 | https://en.wikipedia.org/wiki/Hospice | Hospice | Hospice care is a type of health care that focuses on the palliation of a terminally ill patient's pain and symptoms and attending to their emotional and spiritual needs at the end of life. Hospice care prioritizes comfort and quality of life by reducing pain and suffering. Hospice care provides an alternative to therapies focused on life-prolonging measures that may be arduous, likely to cause more symptoms, or are not aligned with a person's goals.
Hospice care in the United States is largely defined by the practices of the Medicare system and other health insurance providers, which cover inpatient or at-home hospice care for patients with terminal diseases who are estimated to live six months or less. Hospice care under the Medicare Hospice Benefit requires documentation from two physicians estimating a person has less than six months to live if the disease follows its usual course. Hospice benefits include access to a multidisciplinary treatment team specialized in end-of-life care and can be accessed in the home, long-term care facility or the hospital.
Outside the United States, the term tends to be primarily associated with the particular buildings or institutions that specialize in such care. Such institutions may similarly provide care mostly in an end-of-life setting, but they may also be available for patients with other palliative care needs. Hospice care includes assistance for patients' families to help them cope with what is happening and provide care and support to keep the patient at home.
The English word hospice is a borrowing from French. In France however, the word refers more generally to an institution where sick and destitute people are cared for, and does not necessarily have a palliative connotation.
Philosophy
The goal of hospice care is to prioritize comfort, quality of life and individual wishes. How comfort is defined is up to each individual or, if the patient is incapacitated, the patient's family. This can include addressing physical, emotional, spiritual and/or social needs. In hospice care, patient-directed goals are integral and interwoven throughout the care. Hospices typically do not perform treatments that are meant to diagnose or cure an illness but also do not include treatments that hasten death. Instead, hospices focus on palliative care to relieve pain and symptoms.
This philosophy affects how hospice staff treat people and their families. Compared to general healthcare providers, hospice professionals take a different approach to talking to people and their families. They are more likely to make predictions or express uncertainty around future events (e.g., "He might die this week" or "I think she might live longer") than to issue orders or prescribe actions (e.g., "She needs a nurse" or "He can't go home").
History overview
Early development
The word hospice derives from Latin , meaning hospitality or place of rest and protection for the ill and weary. Historians believe the first hospices originated in Malta around 1065, dedicated to caring for the ill and dying en route to and from the Holy Land. The rise of the European Crusading movement in the 1090s placed the incurably ill into places dedicated to treatment. In the early 14th century, the order of the Knights Hospitaller of St. John of Jerusalem opened the first hospice in Rhodes. Hospices flourished in the Middle Ages, but languished as religious orders became dispersed. They were revived in the 17th century in France by the Daughters of Charity of Saint Vincent de Paul. France continued to see development in the hospice field; the hospice of L'Association des Dames du Calvaire, founded by Jeanne Garnier, opened in 1843. Six other hospices followed before 1900.
Meanwhile, hospices developed in other areas. In the United Kingdom attention was drawn to the needs of the terminally ill in the middle of the 19th century, with Lancet and the British Medical Journal publishing articles pointing to the need of the impoverished terminally ill for good care and sanitary conditions. Steps were taken to remedy inadequate facilities with the opening of the Friedenheim in London, which by 1892 offered 35 beds to patients dying of tuberculosis. Four more hospices were established in London by 1905, including the Hostel of God on Clapham Common founded in 1891 by Clara Maria Hole, Mother Superior of Sisterhood of St James' (Anglican) and taken over in 1896 by the Society of Saint Margaret of East Grinstead. Australia, too, saw active hospice development, with notable hospices including the Home for Incurables in Adelaide (1879), the Home of Peace (1902) and the Anglican House of Peace for the Dying in Sydney (1907). In 1899 New York City, the Servants for Relief of Incurable Cancer opened St. Rose's Hospice, which soon expanded to six locations in other cities.
The more influential early developers of hospice included the Irish Religious Sisters of Charity, who opened Our Lady's Hospice in Harold's Cross, Dublin, Ireland, in 1879. It served as many as 20,000 people—primarily with tuberculosis and cancer—dying there between 1845 and 1945. The Sisters of Charity expanded internationally, opening the Sacred Heart Hospice for the Dying in Sydney in 1890, with hospices in Melbourne and New South Wales following in the 1930s. In 1905, they opened St Joseph's Hospice in London.
Hospice movement
In Western society, the concept of hospice began evolving in Europe in the 11th century. In Roman Catholic tradition, hospices were places of hospitality for the sick, wounded, or dying, as well as for travelers and pilgrims. The modern hospice concept includes palliative care for the incurably ill in institutions as hospitals and nursing homes, along with at-home care. The first modern hospice care was created by Dame Cicely Saunders in 1967. Saunders was a British registered nurse whose chronic health problems forced her to pursue a career in medical social work. The relationship she developed with a dying Polish refugee helped solidify her ideas that terminally ill patients needed compassionate care to help address their fears and concerns as well as palliative comfort for physical symptoms. After the refugee's death, Saunders began volunteering at St Luke's Home for the Dying Poor, where a physician told her that she could best influence the treatment of the terminally ill as a physician. Saunders entered medical school while continuing her volunteer work at St. Joseph's. When she completed her degree in 1957, she took a position there.
Saunders emphasized focusing on the patient rather than the disease and introduced the notion of 'total pain', which included psychological and spiritual as well as physical discomfort. She experimented with opioids for controlling physical pain. She also considered the needs of the patient's family. She developed many foundational principles of modern hospice care at St Joseph's.
She disseminated her philosophy internationally in a series of tours of the United States that began in 1963. In 1967, Saunders opened St Christopher's Hospice. Florence Wald, the dean of Yale School of Nursing, who had heard Saunders speak in America, spent a month working with Saunders there in 1969 before bringing the principles of modern hospice care back to the United States, establishing Hospice, Inc. in 1971. Another early hospice program in the United States, Alive Hospice, was founded in Nashville, Tennessee, on November 14, 1975. By 1977 the National Hospice Organization had been formed, and by 1979, a president, Ann G. Blues, had been elected and principles of hospice care had been addressed. At about the same time that Saunders was disseminating her theories and developing her hospice, in 1965, Swiss psychiatrist Elisabeth Kübler-Ross began to consider social responses to terminal illness, which she found inadequate at the Chicago hospital where her American physician husband was employed. Her 1969 best-seller, On Death and Dying, influenced the medical profession's response to the terminally ill. Dr. Balfour Mount introduced the concept of palliative care to Canada in the early 1970s and established the first hospice program at the Royal Victoria Hospital in Montreal, laying the foundation for modern palliative care practices. Saunders and other thanatology pioneers helped to focus attention on the types of care available to them.
In 1984, Josefina Magno, who had been instrumental in forming the American Academy of Hospice and Palliative Medicine and sat as first executive director of the US National Hospice Organization, founded the International Hospice Institute, which in 1996 became the International Hospice Institute and College and later the International Association for Hospice and Palliative Care (IAHPC). The IAHPC follows the philosophy that each country should develop a palliative care model based on its own resources and conditions. IAHPC founding member Derek Doyle told the British Medical Journal in 2003 that Magno had seen "more than 8000 hospice and palliative services established in more than 100 countries." Standards for Palliative and Hospice Care have been developed in countries including Australia, Canada, Hungary, Italy, Japan, Moldova, Norway, Poland, Romania, Spain, Switzerland, the United Kingdom and the United States.
In 2006, the United States–based National Hospice and Palliative Care Organization (NHPCO) and the United Kingdom's Help the Hospices jointly commissioned an independent, international study of worldwide palliative care practices. Their survey found that 15% of the world's countries offered widespread palliative care services with integration into major health care institutions, while an additional 35% offered some form of palliative care services, in some cases localized or limited. As of 2009, an estimated 10,000 programs internationally provided palliative care, although the term hospice is not always employed to describe such services.
In hospice care, the main guardians are the family care giver(s) and a hospice nurse/team who make periodic visits. Hospice can be administered in a nursing home, hospice building, or sometimes a hospital; however, it is most commonly practiced in the home. Hospice care targets the terminally ill who are expected to die within six months.
Popular media
Hospice was the subject of the Netflix 2018 Academy Award–nominated short documentary End Game, about terminally ill patients in a San Francisco hospital and Zen Hospice Project, featuring the work of palliative care physician BJ Miller and other palliative care clinicians. The film was executive produced by hospice and palliative care activist Shoshana R. Ungerleider.
In 2016, an open letter to the singer David Bowie written by a palliative care doctor, Professor Mark Taubert, talked about the importance of good palliative care and hospice provision, especially being able to express wishes about the last months of life, and good education about end of life care generally. The letter went viral after David Bowie's son Duncan Jones shared it. The letter was subsequently read out by the actor Benedict Cumberbatch and the singer Jarvis Cocker at public events.
National variations
Hospice faced resistance from cultural and professional taboos against open communication about death among healthcare providers and the wider population, discomfort with unfamiliar medical techniques and perceived professional callousness towards the terminally ill. Nevertheless, the movement has spread throughout the world.
Africa
A hospice opened in 1980 in Harare (Salisbury), Zimbabwe, the first in Sub-Saharan Africa. In spite of skepticism in the medical community, the hospice movement spread, and in 1987 the Hospice Palliative Care Association of South Africa formed. In 1990, Nairobi Hospice opened in Nairobi, Kenya. As of 2006, Kenya, South Africa and Uganda were among 35 countries offering widespread, well-integrated palliative care. Programs adopted the United Kingdom model, but emphasise home-based assistance.
Following the foundation of hospice in Kenya in the early 1990s, palliative care spread throughout the country. Representatives of Nairobi Hospice sit on the committee to develop a Health Sector Strategic Plan for the Ministry of Health and work with the Ministry of Health to help develop palliative care guidelines for cervical cancer. The Government of Kenya supported hospice by donating land to Nairobi Hospice and providing funding to several of its nurses.
In South Africa, hospice services are widespread, focusing on diverse communities (including orphans and homeless) and offered in diverse settings (including in-patient, day care and home care). Over half of hospice patients in South Africa in the 2003–2004 year were diagnosed with AIDS, with the majority of the remaining diagnosed with cancer. Palliative care is supported by the Hospice Palliative Care Association of South Africa and by national programmes partly funded by the President's Emergency Plan for AIDS Relief.
Hospice Africa Uganda (HAU), founded by Anne Merriman, began offering services in 1993 in a two-bedroom house loaned for the purpose by Nsambya Hospital. HAU has since expanded to a base of operations at Makindye, Kampala, with hospice services offered at roadside clinics by Mobile Hospice Mbarara since January 1998. That same year the Little Hospice Hoima opened in June. Hospice care in Uganda is supported by community volunteers and professionals, as Makerere University offers a distance diploma in palliative care. The government of Uganda published a strategic plan for palliative care that permits nurses and clinical officers from HAU to prescribe morphine.
North America
Canada
Canadian physician Balfour Mount, who first coined the term "palliative care", was a pioneer in medical research and in the Canadian hospice movement, which focused primarily on palliative care in a hospital setting. After meeting Kübler-Ross, Mount studied the experiences of the terminally ill at Royal Victoria Hospital, Montreal; the "abysmal inadequacy", as he termed it, that he found prompted him to spend a week with Cicely Saunders at St. Christopher's. Mount decided to adapt Saunders' model for Canada. Given differences in medical funding, he determined that a hospital-based approach would be more affordable, creating a specialized ward at Royal Victoria in January 1975. Canada's official languages include English and French, leading Mount to propose the term "palliative care ward", as the word hospice was already used in France to refer to nursing homes. Hundreds of palliative care programs then followed throughout Canada through the 1970s and 1980s.
However, as of 2004, according to the Canadian Hospice Palliative Care Association (CHPCA), hospice palliative care was only available to 5–15% of Canadians, with government funding declining. At that time, Canadians were increasingly expressing a desire to die at home, but only two of Canada's ten provinces were provided medication cost coverage for home care. Only four of ten identified palliative care as a core health service. At that time, palliative care was not widely taught at nursing schools or universally certified at medical colleges; only 175 specialized palliative care physicians served all of Canada.
United States
Hospice in the United States has grown from a volunteer-led movement to improve care for people dying alone, isolated, or in hospitals, to a significant part of the health care system. In 2010, an estimated 1.581 million patients received hospice services. Hospice is the only Medicare benefit that includes pharmaceuticals, medical equipment, twenty-four-hour/seven-day-a-week access to care, and support for loved ones following a death. Hospice care is covered by Medicaid and most private insurance plans. Most hospice care is delivered at home. Hospice care is available to people in home-like hospice residences, nursing homes, assisted living facilities, veterans' facilities, hospitals and prisons.
Florence Wald, Dean of the Yale School of Nursing, founded one of the first hospices in the United States in New Haven, Connecticut, in 1974. The first hospital-based palliative care consultation service developed in the US was the Wayne State University School of Medicine in 1985 at Detroit Receiving Hospital. The first US-based palliative medicine and hospice service program was started in 1987 by Declan Walsh at the Cleveland Clinic Cancer Center in Cleveland, Ohio. The program evolved into The Harry R. Horvitz Center for Palliative Medicine, which was designated as a World Health Organization international demonstration project and accredited by the European Society of Medical Oncology as an Integrated Center of Oncology and Palliative Care. Other programs followed; some notable ones are: the Palliative Care Program at the Medical College of Wisconsin (1993); Pain and Palliative Care Service, Memorial Sloan-Kettering Cancer Center (1996); and The Lilian and Benjamin Hertzberg Palliative Care Institute, Mount Sinai School of Medicine (1997).
In 1982, Congress initiated the creation of the Medicare Hospice Benefit, which became permanent in 1986. In 1993, President Clinton installed hospice as a guaranteed benefit and an accepted component of health care provisions. , 1.49 million Medicare beneficiaries were enrolled in hospice care for one day or more, which is a 4.5% increase from the previous year. From 2014 to 2019, Asian- and Hispanic-identifying beneficiaries of hospice care increased by 32% and 21% respectively.
United Kingdom
The first hospice to open in the United Kingdom was the Trinity Hospice in Clapham south London in 1891, on the initiative of the Hoare banking family. More than half a century later, a hospice movement developed after Dame Cicely Saunders opened St Christopher's Hospice in 1967, widely considered the first modern hospice. According to the UK's Help the Hospices, in 2011 UK hospice services consisted of 220 inpatient units for adults with 3,175 beds, 42 inpatient units for children with 334 beds, 288 home care services, 127 hospice at-home services, 272 day care services, and 343 hospital support services. These services together helped over 250,000 patients in 2003 and 2004. Funding varies from 100% funding by the National Health Service to almost 100% funding by charities, but the service is always free to patients. The UK's palliative care has been ranked as the best in the world "due to comprehensive national policies, the extensive integration of palliative care into the National Health Service, a strong hospice movement, and deep community engagement on the issue."
As of 2006, about 4% of all deaths in England and Wales occurred in a hospice setting (about 20,000 patients); a further number of patients spent time in a hospice, or were helped by hospice-based support services, but died elsewhere.
Hospices also provide volunteering opportunities for over 100,000 people in the UK, whose economic value to the hospice movement has been estimated at over £112 million.
Egypt
According to the Global Atlas of Palliative Care at the End of Life, 78% of adults and 98% of children in need of palliative care at the end of life live in low and middle-income countries. Nevertheless, hospice and palliative care provision in Egypt is limited and sparsely available relative to the size of the population. Some of the obstacles to the development of these services have included the lack of public awareness, restricted availability of opioids, and the absence of a national hospice and palliative care development plan. Key efforts made in the past 10 years have been initiated by individuals allowing for the emergence of the first non-governmental organisation providing primarily home-based hospice services in 2010, the opening of one palliative medicine unit at Cairo University in 2008 and an inpatient palliative care unit in Alexandria.
Models of both home-based care and stand-alone hospices exist globally, but with the cultural and societal preferences of patients and their families to die at home in Egypt there is an inclination to focus on the development of home-based hospice and palliative care services.
Israel
The first hospice unit in Israel opened in 1983. More than two decades later, a 2016 study found that 46% of the general Israeli public had never heard of it, despite the 70% of physicians who reported that they had the skill to treat patients according to palliative principles.
Other nations
Hospice care in Australia predated the opening of St Christophers in London by 79 years. The Irish Sisters of Charity opened hospices in Sydney (1889) and in Melbourne (1938). The first hospice in New Zealand opened in 1979. Hospice care entered Poland in the mid-1970s. Japan opened its first hospice in 1981, officially hosting 160 by July 2006. India's first hospice, Shanti Avedna Ashram, opened in Bombay in 1986. The first hospice in the Nordics opened in Tampere, Finland in 1988. The first modern free-standing hospice in China opened in Shanghai in 1988. The first hospice unit in Taiwan, where the term for hospice translates as "peaceful care", opened in 1990. The first free-standing hospice in Hong Kong, where the term for hospice translates as "well-ending service", opened in 1992.
The International Hospice Institute was founded in 1984.
World Hospice and Palliative Care Day
In 2006, the first World Hospice and Palliative Care Day was organised by the Worldwide Palliative Care Alliance, a network of hospice and palliative care national and regional organisations that support the development of hospice and palliative care worldwide. The event takes place on the second Saturday of October every year.
Hospice home health
Nurses that work in hospice in the home healthcare setting aim to relieve pain and holistically support their patient and the patient's family. Patients can receive hospice care when they have less than six months to live or would like to shift the focus of care from curative to comfort care. The goal of hospice care is to meet the needs of both the patient and family, knowing that a home death is not always the best outcome. Medicare covers all costs of hospice treatment.
The hospice home health nurse must be skilled in both physical care and psychosocial care. Most nurses will work with a team that includes a physician, social worker and possibly a spiritual care counselor. Some of the nurse's duties will include reassuring family members, and ensuring adequate pain control. The nurse will need to explain to the patient and family that a pain-free death is possible, and scheduled opioid pain medications are appropriate in this case. The nurse will need to work closely with the medical provider to ensure that dosing is appropriate, and in the case of tolerance, the dose is raised. The nurse should be aware of cultural differences and needs and should aim to meet them. The nurse will also support the family after death and connect the family to bereavement services.
| Biology and health sciences | Medical procedures: General | Health |
27716891 | https://en.wikipedia.org/wiki/Morphism | Morphism | In mathematics, a morphism is a concept of category theory that generalizes structure-preserving maps such as homomorphism between algebraic structures, functions from a set to another set, and continuous functions between topological spaces. Although many examples of morphisms are structure-preserving maps, morphisms need not to be maps, but they can be composed in a way that is similar to function composition.
Morphisms and objects are constituents of a category. Morphisms, also called maps or arrows, relate two objects called the source and the target of the morphism. There is a partial operation, called composition, on the morphisms of a category that is defined if the target of the first object equals the source of the second object. The composition of morphisms behave like function composition (associativity of composition when it is defined, and existence of an identity morphism for every object).
Morphisms and categories recur in much of contemporary mathematics. Originally, they were introduced for homological algebra and algebraic topology. They belong to the foundational tools of Grothendieck's scheme theory, a generalization of algebraic geometry that applies also to algebraic number theory.
Definition
A category C consists of two classes, one of and the other of . There are two objects that are associated to every morphism, the and the . A morphism f from X to Y is a morphism with source X and target Y; it is commonly written as or the latter form being better suited for commutative diagrams.
For many common categories, objects are sets (often with some additional structure) and morphisms are functions from an object to another object. Therefore, the source and the target of a morphism are often called and respectively.
Morphisms are equipped with a partial binary operation, called . The composition of two morphisms f and g is defined precisely when the target of f is the source of g, and is denoted (or sometimes simply gf). The source of is the source of f, and the target of is the target of g. The composition satisfies two axioms:
For every object X, there exists a morphism called the identity morphism on X, such that for every morphism we have .
Associativity whenever all the compositions are defined, i.e. when the target of f is the source of g, and the target of g is the source of h.
For a concrete category (a category in which the objects are sets, possibly with additional structure, and the morphisms are structure-preserving functions), the identity morphism is just the identity function, and composition is just ordinary composition of functions.
The composition of morphisms is often represented by a commutative diagram. For example,
The collection of all morphisms from X to Y is denoted or simply and called the hom-set between X and Y. Some authors write , or . The term hom-set is something of a misnomer, as the collection of morphisms is not required to be a set; a category where is a set for all objects X and Y is called locally small. Because hom-sets may not be sets, some people prefer to use the term "hom-class".
The domain and codomain are in fact part of the information determining a morphism. For example, in the category of sets, where morphisms are functions, two functions may be identical as sets of ordered pairs (may have the same range), while having different codomains. The two functions are distinct from the viewpoint of category theory. Thus many authors require that the hom-classes be disjoint. In practice, this is not a problem because if this disjointness does not hold, it can be assured by appending the domain and codomain to the morphisms (say, as the second and third components of an ordered triple).
Some special morphisms
Monomorphisms and epimorphisms
A morphism is called a monomorphism if implies for all morphisms g1, . A monomorphism can be called a mono for short, and we can use monic as an adjective. A morphism f has a left inverse or is a split monomorphism if there is a morphism such that . Thus is idempotent; that is, . The left inverse g is also called a retraction of f.
Morphisms with left inverses are always monomorphisms, but the converse is not true in general; a monomorphism may fail to have a left inverse. In concrete categories, a function that has a left inverse is injective. Thus in concrete categories, monomorphisms are often, but not always, injective. The condition of being an injection is stronger than that of being a monomorphism, but weaker than that of being a split monomorphism.
Dually to monomorphisms, a morphism is called an epimorphism if implies for all morphisms g1, . An epimorphism can be called an epi for short, and we can use epic as an adjective. A morphism f has a right inverse or is a split epimorphism if there is a morphism such that . The right inverse g is also called a section of f. Morphisms having a right inverse are always epimorphisms, but the converse is not true in general, as an epimorphism may fail to have a right inverse.
If a monomorphism f splits with left inverse g, then g is a split epimorphism with right inverse f. In concrete categories, a function that has a right inverse is surjective. Thus in concrete categories, epimorphisms are often, but not always, surjective. The condition of being a surjection is stronger than that of being an epimorphism, but weaker than that of being a split epimorphism. In the category of sets, the statement that every surjection has a section is equivalent to the axiom of choice.
A morphism that is both an epimorphism and a monomorphism is called a bimorphism.
Isomorphisms
A morphism is called an isomorphism if there exists a morphism such that and . If a morphism has both left-inverse and right-inverse, then the two inverses are equal, so f is an isomorphism, and g is called simply the inverse of f. Inverse morphisms, if they exist, are unique. The inverse g is also an isomorphism, with inverse f. Two objects with an isomorphism between them are said to be isomorphic or equivalent.
While every isomorphism is a bimorphism, a bimorphism is not necessarily an isomorphism. For example, in the category of commutative rings the inclusion is a bimorphism that is not an isomorphism. However, any morphism that is both an epimorphism and a split monomorphism, or both a monomorphism and a split epimorphism, must be an isomorphism. A category, such as a Set, in which every bimorphism is an isomorphism is known as a balanced category.
Endomorphisms and automorphisms
A morphism (that is, a morphism with identical source and target) is an endomorphism of X. A split endomorphism is an idempotent endomorphism f if f admits a decomposition with . In particular, the Karoubi envelope of a category splits every idempotent morphism.
An automorphism is a morphism that is both an endomorphism and an isomorphism. In every category, the automorphisms of an object always form a group, called the automorphism group of the object.
Examples
For algebraic structures commonly considered in algebra, such as groups, rings, modules, etc., the morphisms are usually the homomorphisms, and the notions of isomorphism, automorphism, endomorphism, epimorphism, and monomorphism are the same as the above defined ones. However, in the case of rings, "epimorphism" is often considered as a synonym of "surjection", although there are ring epimorphisms that are not surjective (e.g., when embedding the integers in the rational numbers).
In the category of topological spaces, the morphisms are the continuous functions and isomorphisms are called homeomorphisms. There are bijections (that is, isomorphisms of sets) that are not homeomorphisms.
In the category of smooth manifolds, the morphisms are the smooth functions and isomorphisms are called diffeomorphisms.
In the category of small categories, the morphisms are functors.
In a functor category, the morphisms are natural transformations.
For more examples, see Category theory.
| Mathematics | Category theory | null |
927051 | https://en.wikipedia.org/wiki/RYB%20color%20model | RYB color model | RYB (an abbreviation of red–yellow–blue) is a subtractive color model used in art and applied design in which red, yellow, and blue pigments are considered primary colors. Under traditional color theory, this set of primary colors was advocated by Moses Harris, Michel Eugène Chevreul, Johannes Itten and Josef Albers, and applied by countless artists and designers. The RYB color model underpinned the color curriculum of the Bauhaus, Ulm School of Design and numerous art and design schools that were influenced by the Bauhaus, including the IIT Institute of Design (founded as the New Bauhaus), Black Mountain College, Design Department Yale University, the Shillito Design School, Sydney, and Parsons School of Design, New York.
In this context, the term primary color refers to three exemplar colors (red, yellow, and blue) as opposed to specific pigments. As illustrated, in the RYB color model, red, yellow, and blue are intermixed to create secondary color segments of orange, green, and purple. This set of primary colors emerged at a time when access to a large range of pigments was limited by availability and cost, and it encouraged artists and designers to explore the many nuances of color through mixing and intermixing a limited range of pigment colors. In art and design education, gray, red, yellow, and blue pigments were usually augmented with white and black pigments, enabling the creation of a larger gamut of color nuances including tints and shades.
Although scientifically obsolete because it does not meet the definition of a complementary color in which a neutral or black color must be mixed, it is still a model used in artistic environments, causing confusion about primary and complementary colors. It can be considered an approximation of the CMY color model.
The RYB color model relates specifically to color in the form of paint and pigment application in art and design. Other common color models include the light model (RGB) and the paint, pigment and ink CMY color model, which is much more accurate in terms of color gamut and intensity compared to the traditional RYB color model, the latter emerging in conjunction with the CMYK color model in the printing industry.
History
The first scholars to propose that there are three primary colors for painters were Scarmiglioni (1601), Savot (1609), de Boodt (1609) and Aguilonius (1613). From these, the most influential was the work of Franciscus Aguilonius (1567–1617), although he did not arrange the colors in a wheel.
Jacob Christoph Le Blon was the first to apply the RYB color model to printing, specifically mezzotint printing, and he used separate plates for each color: yellow, red and blue plus black to add shades and contrast. In 'Coloritto', Le Blon asserted that “the art of mixing colours…(in) painting can represent all visible objects with three colours: yellow, red and blue; for all colours can be composed of these three, which I call Primitive”. Le Blon added that red and yellow make orange; red and blue, make purple; and blue and yellow make green (Le Blon, 1725, p6).
In the 18th century, Moses Harris advocated that a multitude of colors can be created from three "primitive" colors – red, yellow, and blue.
Mérimée referred to "three simple colours (yellow, red, and blue)" that can produce a large gamut of color nuances. "United in pairs, these three primitive colours give birth to three other colours as distinct and brilliant as their originals; thus, yellow mixed with red, gives orange; red and blue, violet; and green is obtained by mixing blue and yellow" (Mérimée, 1839, p245). Mérimée illustrated these color relationships with a simple diagram located between pages 244 and 245: Chromatic Scale (Echelle Chromatique).De la peinture à l’huile : ou, Des procédés matériels employés dans ce genre de peinture, depuis Hubert et Jean Van-Eyck jusqu’à nos jours was published in 1830 and an English translation by W. B. Sarsfield Taylor was published in London in 1839.
Similar ideas about the creation of color using red, yellow, and blue were discussed in Theory of Colours (1810) by the German poet, color theorist and government minister Johann Wolfgang von Goethe.
In The Law of Simultaneous Color Contrast (1839) by the French industrial chemist Michel Eugène Chevreul discussed the creation of numerous color nuances and his color theories were underpinned by the RYB color model.
Separate to the RYB color model, cyan, magenta, and yellow primary colors are associated with CMYK commonly used in the printing industry. Cyan, magenta, and yellow are often referred to as "process blue", "process red", and "process yellow".
Old model of coloration with four primaries
The ancient Greeks, under the influence of Aristotle, Democritus and Plato, considered that there were four basic colors that coincided with the four elements: earth (ochre), sky (blue), water (green) and fire (red), while black and white represented the light of day and the darkness of night. The four-color system is formed by the primaries yellow, green, blue and red, and was supported by Alberti in his "De Pictura" (1436), using the rectangle, rhombus, and color wheel to represent them.
Leonardo da Vinci endorsed this model in 1510, although he hesitated to include green, noting that green could be obtained by mixing blue and yellow. Also Richard Waller, in his "Catalogue of Simple and Mixed Colors" (1686), graphed these four colors in a square. These four colors have often been referred to as "the primary psychological colors".
Traditional coloring with three primaries
The first known case of trichromacy coloration (of 3 primaries) can be found in a work on optics by the Belgian thinker Franciscus Aguilonius in 1613, who in his "Opticorum libri sex, philosophis iuxtà ac mathematicis utiles" in Latin (Roughly, Six books of optics: useful to philosophers as well as to mathematicians), graphed the colors flavvus, rvbevs and cærvlevs (yellow, red and blue) giving rise to the intermediate colors avrevs, viridis and pvrpvrevs (orange, green and purple) and their relationship with the extremes albvs and niger (white and black). However, the idea of three primary colors is older, as Aguilonius supported the view known since the Middle Ages that the colors yellow, red, and blue were the basic or "noble" colors from which all others are derived.
This model was used for printing by Jacob Christoph Le Blon in 1725 and called it Coloritto or harmony of colouring, stating that the primitive (primary) colors are yellow, red and blue, while the secondary are orange, green and purple or violet.
In 1766, Moses Harris developed an 18-color color wheel based on this model, including a wider range of colors by adding light and dark derivatives. During the 18th and 19th centuries, this color model was endorsed by many authors who have left illustrations that can still be appreciated today, such as Louis-Bertrand Castel (1740), the Tobias's color system Mayer (1758), Moses Harris (1770–76), Ignaz Schiffermuller (1772), Baumgartner and Muller (1803), Sowerby (1809), Runge (1809), the popular "Theory of Colors" (1810) by Goethe, Gregoire (1810–20), Merimee (1815-30-39), Klotz (1816), G. Field (1817-41-50), Hayter (1826 ), the "Law of Simultaneous Contrast of Colours" (1839) by Chevreul and many others.
By the 20th century, natural pigments gave way to synthetic ones. The invention of phthalocyanine and derivatives of quinacridone, expanded the range of primary blues and reds, getting closer to the ideal subtractive colors and the CMY and CMYK models.
| Physical sciences | Basics | Physics |
927572 | https://en.wikipedia.org/wiki/Guernsey%20cattle | Guernsey cattle | The Guernsey is a breed of dairy cattle from the island of Guernsey in the Channel Islands. It is fawn or red and white in colour, and is hardy and docile. Its milk is rich in flavour, high in fat and protein, and has a golden-yellow tinge due to its high β-carotene content. The Guernsey is one of three Channel Island cattle breeds; the other two are the Alderney, which is now extinct, and the Jersey.
History
The Guernsey was bred on the Channel Island of Guernsey; it is first documented in the nineteenth century, and its origins are unknown. Cattle were brought to the island in the Middle Ages for draught work. It has been suggested that the Guernsey derives from cattle imported from the French mainland – brindled cattle from Normandy, and wheaten stock similar to the Froment du Léon of Brittany. There may also have been some influence from Dutch cattle in the 18th century. During that century large numbers of cattle were exported from the Channel Islands to England; some of them had previously been brought from France. Imports of French cattle to Guernsey were forbidden by law in 1819, but some importation of British cattle continued until 1877. Some cattle evacuated from Alderney during the Second World War were merged into the breed.
Exports of cattle and semen were for a while an important economic resource for the island, and in the early 20th century, a large number of Guernsey cattle were exported to the United States. The Guernsey breed is on the watch list maintained by the American Livestock Breeds Conservancy, with fewer than 2,500 annual registrations in the U.S. and an estimated global population less than 10,000 animals.
Characteristics
The Guernsey is of medium size: cows weigh , and bulls . The coat is red or fawn (wheat-coloured), and may or may not be pied red-and-white or fawn-and-white. The Guernsey produces rich and flavoursome milk. It traditionally had several other good qualities: it was long-lived, calved without difficulty, grazed well and – being relatively small-sized – was an efficient milk producer. These advantages have been compromised by recent selective breeding strategies, which have led to larger animals, with longer legs. These no longer display the traditional qualities of the breed; this is particularly marked where there has been cross-breeding with Holstein-Friesian stock.
Use
The Guernsey is a dairy breed, and generally is reared for that purpose only. However, the cow is usually removed from the dairy herd around ages six to eight, and marketed for beef, and other processed meats. The milk has a golden-yellow tinge due to a high content of β-carotene, a provitamin for vitamin A. The milk also has a high butterfat content of 5% and a high protein content of 3.7%. Guernsey cows produce around 6000 litres per cow per year.
| Biology and health sciences | Cattle | Animals |
927752 | https://en.wikipedia.org/wiki/HFS%20Plus | HFS Plus | HFS Plus or HFS+ (also known as Mac OS Extended or HFS Extended) is a journaling file system developed by Apple Inc. It replaced the Hierarchical File System (HFS) as the primary file system of Apple computers with the 1998 release of Mac OS 8.1. HFS+ continued as the primary Mac OS X file system until it was itself replaced with the Apple File System (APFS), released with macOS High Sierra in 2017. HFS+ is also one of the formats supported by the iPod digital music player.
Compared to its predecessor HFS, also called Mac OS Standard or HFS Standard, HFS Plus supports much larger files (block addresses are 32-bit length instead of 16-bit) and using Unicode (instead of Mac OS Roman or any of several other character sets) for naming items. Like HFS, HFS Plus uses B-trees to store most volume metadata, but unlike most file systems that support hard links, HFS Plus supports hard links to directories. HFS Plus permits filenames up to 255 characters in length, and n-forked files similar to NTFS, though until 2005 almost no system software took advantage of forks other than the data fork and resource fork. HFS Plus also uses a full 32-bit allocation mapping table rather than HFS's 16 bits, improving the use of space on large disks.
History
Codenamed Sequoia in development, HFS+ was introduced with the January 19, 1998, release of Mac OS 8.1.
With the release of the Mac OS X 10.2.2 update on November 11, 2002, Apple added optional journaling features to HFS Plus for improved data reliability. These features were accessible through the GUI, using the Disk Utility application in Mac OS X Server, but only accessible through the command line in the standard desktop client.
With Mac OS X v10.3, all HFS Plus volumes on all Macs were set to be journaled by default. Within the system, an HFS Plus volume with a journal is identified as HFSJ.
Mac OS X 10.3 also introduced another version of HFS Plus called HFSX. HFSX volumes are almost identical to HFS Plus volumes, except that they are never surrounded by the HFS Wrapper that is typical of HFS Plus volumes and they optionally support case sensitivity for file and folder names. HFSX volumes can be recognized by two entries in the Volume Header, a value of HX in the signature field and 5 in the version field.
Mac OS X 10.3 also marked Apple's adoption of Unicode 3.2 decomposition, superseding the Unicode 2.1 decomposition used previously. This change caused problems for developers writing software for Mac OS X.
Mac OS X 10.3 introduced a number of techniques that are intended to avoid fragmentating files in HFS+.
With Mac OS X 10.4, Apple added support for Inline Attribute Data records, something that had been a part of the Mac OS X implementation of HFS Plus since at least 10.0, but always marked as "reserved for future use". Until the release of Mac OS X Server 10.4, HFS Plus supported only the standard UNIX file system permissions; however, 10.4 introduced support for access control list–based file security, which provides a richer mechanism to define file permissions and is also designed to be fully compatible with the file permission models on other platforms such as Microsoft Windows XP and Windows Server 2003.
In Mac OS X Leopard 10.5, directory hard-linking was added as a fundamental part of Time Machine.
In Mac OS X Snow Leopard 10.6, HFS+ compression was added using Deflate (Zlib). In open source and some other areas this is referred to as AppleFSCompression or decmpfs. Compressed data may be stored in either an extended attribute or the resource fork. When using non-Apple APIs, AppleFSCompression is not always completely transparent. OS X 10.9 introduced two new algorithms: LZVN (libFastCompression), and LZFSE.
In Mac OS X Lion 10.7, logical volume encryption (known as FileVault 2) was added to the operating system. This addition to the operating system in no way changed the logical structure of the file system. Apple's logical volume manager is known as Core Storage and its encryption at the volume level can apply to file systems other than HFS Plus. With appropriate hardware, both encryption and decryption should be transparent.
Design
HFS Plus volumes are divided into sectors (called logical blocks in HFS), that are usually 512 bytes in size. These sectors are then grouped together into allocation blocks which can contain one or more sectors; the number of allocation blocks depends on the total size of the volume. HFS Plus uses a larger value to address allocation blocks than HFS, 32 bits rather than 16 bits; this means it can access 4,294,967,296 (= 232) allocation blocks rather than the 65,536 (= 216) allocation blocks available to HFS. When disks were small, this was of little consequence, but as larger-capacity drives became available, it meant that the smallest amount of space that any file could occupy (a single allocation block) became excessively large, wasting significant amounts of space. For example, on a 1 GB disk, the allocation block size under HFS is 16 KB, so even a 1-byte file would take up 16 KB of disk space. HFS Plus's system greatly improves space utilization on larger disks as a result.
File and folder names in HFS Plus are also encoded in UTF-16 and normalized to a form very nearly the same as Unicode Normalization Form D (NFD) (which means that precomposed characters like "å" are decomposed in the HFS+ filename and therefore count as two code units and UTF-16 implies that characters from outside the Basic Multilingual Plane also count as two code units in an HFS+ filename). HFS Plus permits filenames up to 255 UTF-16 code units in length.
Formerly, HFS Plus volumes were embedded inside an HFS standard file system. This was phased out by the Tiger transition to Intel Macs, where the HFS Plus file system was not embedded inside a wrapper. The wrapper had been designed for two purposes; it allowed Macintosh computers without HFS Plus support in their ROM to boot HFS Plus volumes and it also was designed to help users transition to HFS Plus by including a minimal HFS volume with a read-only file called Where_have_all_my_files_gone?, explaining to users with versions of Mac OS 8.0 and earlier without HFS Plus, that the volume requires a system with HFS Plus support. The original HFS volume contains a signature and an offset to the embedded HFS Plus volume within its volume header. All allocation blocks in the HFS volume which contain the embedded volume are mapped out of the HFS allocation file as bad blocks.
Notable among file systems used for Unix systems, HFS Plus does not support sparse files.
There are nine structures that make up a typical HFS Plus volume:
Sectors 0 and 1 of the volume are HFS boot blocks. These are identical to the boot blocks in an HFS volume. They are part of the HFS wrapper.
Sector 2 contains the Volume Header, which is equivalent to the Master Directory Block in an HFS volume. The Volume Header stores a wide variety of data about the volume itself, for example the size of allocation blocks, a timestamp that indicates when the volume was created or the location of other volume structures such as the Catalog File or Extent Overflow File. The Volume Header is always located in the same place.
The Allocation File which keeps track of which allocation blocks are free and which are in use. It is similar to the Volume Bitmap in HFS, in which each allocation block is represented by one bit. A zero means the block is free and a one means the block is in use. The main difference with the HFS Volume Bitmap, is that the Allocation File is stored as a regular file – it does not occupy a special reserved space near the beginning of the volume. The Allocation File can also change size and does not have to be stored contiguously within a volume.
The Catalog File is a B-tree that contains records for all the files and directories stored in the volume. The HFS Plus Catalog File is very similar to the HFS Catalog File, the main differences being records are larger to allow more fields and to allow for those fields to be larger (for example to allow the longer 255-character unicode file names in HFS Plus). A record in the HFS Catalog File is 512 bytes in size; a record in the HFS Plus Catalog File is 4 KB in the classic Mac OS and 8 KB in macOS. Fields in HFS are of fixed size, while in HFS Plus the size can vary depending on the actual size of the data they store.
The Extents Overflow File is another B-tree that records the allocation blocks that are allocated to each file as extents. Each file record in the Catalog File is capable of recording eight extents for each fork of a file; once those are used additional extents are recorded in the Extents Overflow File. Bad blocks are also recorded as extents in the Extents Overflow File. The default size of an extent record in the classic Mac OS is 1 KB and 4 KB in macOS.
The Attributes File is a new B-tree in HFS Plus that does not have a corresponding structure in HFS. The Attributes File can store three different types of 4 KB records: Inline Data Attribute records, Fork Data Attribute records and Extension Attribute records. Inline Data Attribute records store small attributes that can fit within the record itself. Fork Data Attribute records contain references to a maximum of eight extents that can hold larger attributes. Extension Attributes are used to extend a Fork Data Attribute record when its eight extent records are already used.
The Startup File is designed for non-Mac OS systems that lack HFS or HFS Plus support. It is similar to the Boot Blocks of an HFS volume.
The second-to-last sector contains the Alternate Volume Header, which is equivalent to the Alternate Master Directory Block of HFS. This is the second-to-last-sector for the disk, not the volume; if the disk is larger than the volume, the AVH will be outside the range of the filesystem.
The last sector in the volume is reserved for use by Apple. It is used during the computer manufacturing process.
Criticisms
HFS Plus lacks several features considered staples of modern file systems like ZFS and NTFS. Data checksums are the most routinely cited missing feature.
Besides checksumming, features of modern file systems that HFS+ lacks include:
nanosecond timestamps
concurrent access (that is, more than one process can access the filesystem at the same time)
snapshotting
Support for dates beyond February 6, 2040
sparse file support
a better implementation of hard links (in other filesystems, these are typically multiple directory entries pointing to the same data blocks; hard links in macOS are implemented as small files that are stored in a special hidden directory)
HFS Plus was not designed for Unix-like systems, so features such as file system permissions and hard links had to be retrofitted when Apple moved to Mac OS X.
Other operating systems
Linux
The Linux kernel includes the hfsplus module for mounting HFS+ filesystems read-write. HFS+ fsck and mkfs have been ported to Linux and are part of the hfsprogs package.
In 2009, these drivers were diagnosed to be corrupting HFS+ drives with a capacity greater than 2 TB. Consequently, Linux distributions such as Debian and Ubuntu stopped allowing mounting of HFS+ drives or partitions greater than 2 TB. , work is in progress to lift this restriction.
Under Linux's current HFS+ driver, journaling must be disabled in order to write data safely onto an HFS+ partition. Provided the partition isn't being used by Apple's Time Machine software, journaling can be disabled under macOS: Using Disk Utility in OS X Yosemite, the user may hold Alt/Option and click "Disable Journaling" on the File menu, having first selected a mounted partition.
An HFS+ partition with journaling enabled may be forcibly mounted with write access under Linux, but this is unsupported and unwise.
A Google Summer of Code project to implement write support to journaled HFS+ was accepted by the Linux Foundation in 2011 but was not completed at that time and is still a work in progress. Progress and improvements to the HFS+ driver, including some updates to journaling support, are posted on the linux-fsdevel mailing list from time to time.
, Paragon Software Group provided kernel drivers that allow full read-write access to HFS+ journaled volumes. The product is a proprietary implementation of HFS+ based on Paragon's proprietary UFSD library. There are both free and paid editions of the driver, and they include a utility for checking and repairing HFS+ volumes. According to the online documentation (free version or the paid edition), both the free edition and the paid edition currently support Linux kernels from 2.6.36 up to 4.12.x. Ubuntu, Debian, Fedora Linux, Rocky Linux, Red Hat Enterprise Linux, OpenSUSE and CentOS are the only Linux distributions officially supported.
Windows
As of May 2012, Apple has only released read-only HFS+ drivers for Windows XP, Windows Vista, and Windows 7 as part of the Boot Camp software in Mac OS X 10.6. This means users on these systems can read data on the HFS+ drive, but not write to them. Microsoft has created an HFS+ driver for the Xbox 360 mainly for the purpose of reading HFS+-formatted iPods.
A free and opensource software – jHFSplus, based on HFSExplorer and jpfm – can be used to mount hfs/hfs+ partitions as read-only virtual folders.
A freeware plugin for Total Commander exists, that can read, among others, HFS and HFS+ filesystems.
DiskInternals Linux Reader can be used to extract/save folders/files out of HFS and HFS+ Hard Drives/Partitions.
A commercial product, MacDrive, is also available for mounting HFS and HFS+ drives, optical discs, and other media in Windows Explorer, and allows both reading and writing to the volume, as well as repairing and formatting Mac disks.
A commercial product, Paragon's HFS+ for Windows allows full read and write and disk management from all versions of Windows from Windows XP to Windows Server 2008.
Cross-platform
A free (GPL) alternative is HFSExplorer written by Erik Larsson. HFSExplorer is a Java application for viewing and extracting files from an HFS+ volume (Mac OS Extended) or an HFSX volume (Mac OS Extended, Case-sensitive). The volume can be located either on a physical disk, in various Apple disk image and sparse disk image formats, or a raw file system dump. However, HFSExplorer is a read-only solution; it cannot write to HFS-formatted volumes.
| Technology | Data storage and memory | null |
928680 | https://en.wikipedia.org/wiki/Grass%20snake | Grass snake | The grass snake (Natrix natrix), sometimes called the ringed snake or water snake, is a Eurasian semi-aquatic non-venomous colubrid snake. It is often found near water and feeds almost exclusively on amphibians.
Subspecies
Many subspecies are recognized, including:
Natrix natrix helvetica (Lacépède, 1789) was formerly treated as a subspecies, but following genetic analysis it was recognised in August 2017 as a separate species, Natrix helvetica, the barred grass snake. Four other subspecies were transferred from N. natrix to N. helvetica, becoming N. helvetica cettii, N. helvetica corsa, N. helvetica lanzai and N. helvetica sicula.
Description
The grass snake is typically dark green or brown in colour with a characteristic yellow or whitish collar behind the head, which explains the alternative name ringed snake. The colour may also range from grey to black, with darker colours being more prevalent in colder regions, presumably owing to the thermal benefits of being dark in colour. The underside is whitish with irregular blocks of black, which are useful in recognizing individuals. It can grow to or more in length.
Distribution
The grass snake is widely distributed in mainland Europe, ranging from mid Scandinavia to southern Italy. It is also found in the Middle East and northwestern Africa.
Grass snakes in Britain were thought to belong to the subspecies N. n. helvetica but have been reclassified as the barred grass snake Natrix helvetica. Any records of N. natrix in Britain are now considered to have originated from imported specimens.
Ecology
Feeding
Grass snakes mainly prey on amphibians, especially the common toad and the common frog, although they may also occasionally eat ants and larvae. Captive snakes have been observed taking earthworms offered by hand, but dead prey items are never taken. The snake will search actively for prey, often on the edges of the water, using sight and sense of smell (using Jacobson's organ). They consume prey live without using constriction.
Habitat
Grass snakes are strong swimmers and may be found close to fresh water, although there is evidence individual snakes often do not need bodies of water throughout the entire season.
The preferred habitat appears to be open woodland and "edge" habitat, such as field margins and woodland borders, as these may offer adequate refuge while still affording ample opportunity for thermoregulation through basking. Pond edges are also favoured and the relatively high chance of observing this secretive species in such areas may account for their perceived association with ponds and water.
Grass snakes, like most reptiles, are at the mercy of the thermal environment and need to overwinter in areas which are not subject to freezing. Thus, they typically spend the winter underground where the temperature is relatively stable.
Reproduction
As spring approaches, the males emerge first and spend much of the day basking in an effort to raise body temperature and thereby metabolism. This may be a tactic to maximise sperm production, as the males mate with the females as soon as they emerge up to two weeks later in April, or earlier if environmental temperatures are favourable. The leathery-skinned eggs are laid in batches of eight to 40 in June to July and hatch after about 10 weeks. To survive and hatch, the eggs require a temperature of at least , but preferably , with high humidity. Areas of rotting vegetation, such as compost heaps, are preferred locations. The young are about long when they hatch and are immediately independent.
Migration
After breeding in summer, snakes tend to hunt and may range widely during this time, moving up to several hundred metres in a day. Prey items tend to be large compared to the size of the snake, and this impairs the movement ability of the snake. Snakes that have recently eaten rarely move any significant distance and will stay in one location, basking to optimize their body temperature until the prey item has been digested. Individual snakes may only need two or three significant prey items throughout an entire season.
Ecdysis (moulting)
Ecdysis occurs at least once during the active season. As the outer skin wears and the snake grows, the new skin forms underneath the old, including the eye scales which may turn a milky blue/white colour at this time—referred to as being 'in blue'. The blue-white colour comes from an oily secretion between the old and new skins; the snake's coloration will also look dull, as though the animal is dusty. This process affects the eyesight of the snakes and they do not move or hunt during this time; they are also, in common with most other snakes, more aggressive. The outer skin is eventually sloughed in one piece (inside-out) and normal movement activity is resumed.
Defence
In defence they can produce a garlic-smelling fluid from the anal glands, and feign death (thanatosis) by becoming completely limp when they may also secrete blood (autohaemorrhage) from the mouth and nose. They may also perform an aggressive display in defence, hissing and striking without opening the mouth. They rarely bite in defense and lack venomous fangs. When caught they often regurgitate the contents of their stomachs.
Grass snakes display a rare defensive behavior involving raising the front of the body and flattening the head and neck so that it resembles a cobra's hood, although the geographic ranges of grass snakes and of cobras overlap very little. However, the fossil record shows that the extinct European cobra Naja romani occurs in Miocene-aged strata of France, Germany, Austria, Romania, and Ukraine and thus overlapped with Natrix species including the extinct Natrix longivertebrata, suggesting that the grass snake's behavioral mimicry of cobras is a fossil behavior, although it may protect against predatory birds which migrate to Africa for the winter and encounter cobras there.
Protection and threats
The species has various predator species, including corvids, storks, owls and perhaps other birds of prey, foxes, and the domestic cat.
In Denmark it is protected, as all five species of reptiles were protected in 1981. Two of the subspecies are considered critically endangered: N. n. cetti (Sardinian grass snake) and N. n. schweizeri.
Mythology
Baltic
In the Baltic mythology, the grass snake (Lithuanian: žaltys, Latvian: zalktis) is seen as a sacred animal. It was frequently kept as a pet, living under a married couple's bed or in a special place near the hearth. Supposedly, snakes ate food given to them by hand.
After the Christianization of Lithuania and Latvia, the grass snake still retained some mythological significance. In spite of the serpent's symbolic meaning as a symbol of evil in Christianity, in Latvia and Lithuania there were various folk beliefs, dating even to the late 19th century, that killing grass snakes might bring grave misfortune or that an injured snake will take revenge on the offender. The ancient Baltic belief of grass snakes as household spirits transformed into a belief that there is a snake (known or not to the inhabitants) living under every house; if it leaves, the house will burn down. Common Latvian folk sayings include "who kills a grass snake, kills his happiness" and "when the Saulė sees a dead grass snake, she cries for 9 days".
Well-known literary works based on these traditions include Lithuanian folk tale Eglė the Queen of Serpents (Eglė žalčių karalienė) and the Latvian folk fairytale "The grass snake's bride" (Zalkša līgava). These works include another common theme in Baltic mythology: that grass snakes wear crowns (note grass snake's yellow spots) and that there is a king of snakes who wears a golden crown. In some traditions the king of snakes changes every year; he drops his crown in spring and the other snakes fight for it (possibly based on the mating of grass snakes).
Today grass snakes hold a meaning of house blessing among many Latvians and Lithuanians. One tradition is to put a bowl of milk near a snake's place of residence, although there is no evidence of a grass snake ever drinking milk. Driven by late 19th century and 20th century Romantic nationalism, grass snake motifs in Latvia have gained a meaning of education and wisdom, and are common ornaments in the military, folk dance groups and education logos and insignia. They are also found on the Lielvārde Belt. The grass snake has also become one of the main symbols of the Lithuanian neo-pagan movement Romuva.
Roman
Virgil in his 29 BC Georgics (book III, lines 425-439: ) describes the grass snake as a large feared snake living in marshes in Calabria, eating frogs and fish.
Gallery
| Biology and health sciences | Snakes | Animals |
929532 | https://en.wikipedia.org/wiki/Haemophilus%20influenzae | Haemophilus influenzae | Haemophilus influenzae (formerly called Pfeiffer's bacillus or Bacillus influenzae) is a Gram-negative, non-motile, coccobacillary, facultatively anaerobic, capnophilic pathogenic bacterium of the family Pasteurellaceae. The bacteria are mesophilic and grow best at temperatures between 35 and 37 °C.
H. influenzae was first described in 1893 by Richard Pfeiffer during an influenza pandemic when he incorrectly identified it as the causative microbe, which is why the bacteria was given the name "influenzae". H. influenzae is responsible for a wide range of localized and invasive infections, typically in infants and children, including pneumonia, meningitis, or bloodstream infections. Treatment consists of antibiotics; however, H. influenzae is often resistant to the penicillin family, but amoxicillin/clavulanic acid can be used in mild cases. Serotype B H. influenzae have been a major cause of meningitis in infants and small children, frequently causing deafness and mental retardation. However, the development in the 1980s of a vaccine effective in this age group (the Hib vaccine) has almost eliminated this in developed countries.
This species was the first organism to have its entire genome sequenced.
Physiology and metabolism
Structure
H. influenzae is a small Gram-negative bacterium, approximately 0.3 micrometer to 1 micrometer. Like other Gram-negative bacteria, H. influenzae has a thin peptidoglycan layer surrounded by an outer membrane containing lipopolysaccharide. Some types of H. influenzae contain a polysaccharide capsule around the outer membrane to aid in protection and colonization. The bacteria are pleomorphic, meaning the shape of the bacterium is variable, however it is typically coccobacillus or rod-shaped. H. Influenzae contains pili, which are specialized to adhere to the human nasopharynx. The H. Influenzae pili, unlike those of E. coli, resist unwinding, allowing for stronger adhesion to resist expulsion when coughing or sneezing. A minority of non-typeable, or unencapsulated, H. influenzae employ a variety of attachment techniques, such as pili, adhesins, or Hia and Hap proteins. Though the bacteria possess pili, they are not used for traditional movement or motility, and the bacterium is still considered to be non-motile.
The cell wall of H. influenzae bacterium contains various proteins, referred to as autotransporters, for adherence and colony formation. H. influenzae prefers to bind to mucus linings or non-ciliated epithelial cells, which is facilitated by Hap𝘴 autotransporters in the cell wall binding with unknown receptors within the epithelium. The Hap𝘴 autotransporters also facilitate the formation of microcolonies of the bacteria. These microcolonies are likely responsible for the formation of various biofilms within the body, such as those responsible for middle ear or lung infections.
Penicillin binding proteins
Penicillin binding proteins (PBPs) catalyze steps in peptidoglycan metabolism. They carry out essential processes needed to build and modify the cell wall. These proteins are the targets blocked by penicillin and other beta-lactam antibiotics that bind to PBPs, hence their name. Some antibiotic-resistant isolates of H. Influenzae contain modified PBPs that resist beta-lactam action by producing beta-lactamases to degrade these antibiotics. This resistance is likely due to a N526K mutation, or R517H substitution in conjunction with another unknown mutation. The R517H substitution alone did not have a lower affinity for penicillin, and therefore cannot cause resistance alone. Beta-lactamase emergence in the 1970s caused the therapy for severe cases of H. influenzae to be changed from ampicillin to cephalosporins, however further resistance to cephalosporins has occurred due to changes in the transpeptidase domain of penicillin binding protein 3 (PBP3).
Serotypes
H. influenzae isolates were initially characterized as either encapsulated (having an extracellular polysaccharide layer, the bacterial capsule) or unencapsulated. Encapsulated strains were further classified on the basis of the immune response to the type of polysaccharides in their capsule. The six generally recognized types of encapsulated H. influenzae are: a, b, c, d, e, and f. H. Influenzae type b, also known as Hib, is the most common form, recognizable by its polyribosyl ribitol phosphate (PRP) capsule, and found mostly in children. Types a, e, and f have been isolated infrequently, while types d and c are rarely isolated. Unencapsulated strains are more genetically diverse than the encapsulated group. Unencapsulated strains are termed nontypable (NTHi) because they lack capsular serotypes; however, all H. influenzae isolates can now be classified by multilocus sequence typing and other molecular methods. Most NTHi strains are considered to be part of the normal human flora in the upper and lower respiratory tract, genitals, and conjunctivae (mucous membranes of the eye).
Metabolism
H. influenzae uses the Embden–Meyerhof–Parnas (EMP) pathway for glycolysis and the pentose phosphate pathway, which is anabolic rather than catabolic. The citric acid cycle is incomplete and lacks several enzymes that are found in a fully functioning cycle. The enzymes missing from the TCA cycle are citrate synthase, aconitate hydratase, and isocitrate dehydrogenase. H. influenzae has been found in both aerobic and anaerobic environments, as well as environments with different pH's.
Genome and genetics
H. influenzae was the first free-living organism to have its entire genome sequenced. The sequencing was completed by Craig Venter and his team at the Institute for Genomic Research, now part of the J. Craig Venter Institute. Haemophilus was chosen because one of the project leaders, Nobel laureate Hamilton Smith, had been working on it for decades and was able to provide high-quality DNA libraries. The sequencing method used was whole-genome shotgun, which was completed and published in Science in 1995.
The genome of strain Rd KW20 consists of 1,830,138 base pairs of DNA in a single circular chromosome that contains 1604 protein-coding genes, 117 pseudogenes, 57 tRNA genes, and 23 other RNA genes. About 90% of the genes have homologs in E. coli, another gamma-proteobacterium. In fact, the similarity between genes of the two species ranges from 18% to 98% protein sequence identity, with the majority sharing 40–80% of their amino acids (with an average of 59%).
Conjugative plasmids (DNA molecules that are capable of horizontal transfer between different species of bacteria) can frequently be found in H. influenzae. It is common that the F+ plasmid of a competent Escherichia coli bacterium conjugates into the H. influenzae bacterium, which then allows the plasmid to transfer among H. influenzae strands via conjugation.
Role of transformation
H. influenzae mutants defective in their rec1 gene (a homolog of recA) are very susceptible to being killed by the oxidizing agent hydrogen peroxide. This finding suggests that rec1 expression is important for H. influenzae survival under conditions of oxidative stress. Since it is a homolog of recA, rec1 likely plays a key role in recombinational repair of DNA damage. Thus, H. influenzae may protect its genome against the reactive oxygen species produced by the host's phagocytic cells through recombinational repair of oxidative DNA damages. Recombinational repair of a damaged site of a chromosome requires, in addition to rec1, a second homologous undamaged DNA molecule. Individual H. influenzae cells are capable of taking up homologous DNA from other cells by the process of transformation. Transformation in H. influenzae involves at least 15 gene products, and is likely an adaptation for repairing DNA damage in the resident chromosome.
Culture methods and diagnosis of infections
Culture
Bacterial culture of H. influenzae is performed on agar plates. The strongest growth is seen on chocolate agar at 37 °C in a CO2-enriched incubator. The ideal CO2 concentration for the culture is ~5%. However adequate growth is often seen on brain-heart infusion agar supplemented with hemin and nicotinamide adenine dinucleotide (NAD)
Colonies of H. influenzae appear as convex, smooth, pale, grey, or transparent colonies with a mild odor. H. influenzae will only grow on blood agar if other bacteria are present to release these factors from the red blood cells, forming 'satellite' colonies around these bacteria. For example, H. influenzae will grow in the hemolytic zone of Staphylococcus aureus on blood agar plates; the hemolysis of cells by S. aureus releases NAD which is needed for its growth. H. influenzae will not grow outside the hemolytic zone of S. aureus due to the lack of nutrients in these areas.
Diagnosis of infections
Clinical features of a respiratory tract infection may include initial symptoms of an upper respiratory tract infection mimicking a viral infection, usually associated with low-grade fevers. This may progress to the lower respiratory tract within a few days, with features often resembling those of wheezy bronchitis. Sputum may be difficult to expectorate and is often grey or creamy in color. The cough may persist for weeks without appropriate treatment. Many cases are diagnosed after presenting chest infections that do not respond to penicillins or first-generation cephalosporins. A chest X-ray can identify alveolar consolidation.
Clinical diagnosis of invasive H. influenzae infection (infection that has spread to the bloodstream and internal tissues) is typically confirmed by bacterial culture, latex particle agglutination tests, or polymerase chain reaction tests on clinical samples obtained from an otherwise sterile body site. In this respect, H. influenzae cultured from the nasopharyngeal cavity or throat would not indicate H. influenzae disease, because these sites are colonized in disease-free individuals. However, H. influenzae isolated from cerebrospinal fluid or blood or joint fluid would indicate invasive H. influenzae infection. Microscopic observation of a Gram stained specimen of H. influenzae will show Gram-negative coccobacillus. The cultured organism can be further characterized using catalase and oxidase tests, both of which should be positive. Further serological testing is necessary to distinguish the capsular polysaccharide and differentiate between H. influenzae b and nonencapsulated strains.
Although highly specific, bacterial culture of H. influenzae lacks sensitivity. Use of antibiotics prior to sample collection greatly reduces the isolation rate by killing the bacteria before identification is possible. Recent work has shown that H. influenzae uses a highly specialized spectrum of nutrients where lactate is a preferred carbon source.
Latex particle agglutination
The latex particle agglutination test (LAT) is a more sensitive method to detect H. influenzae than is culture. Because the method relies on antigen rather than viable bacteria, the results are not disrupted by prior antibiotic use. It also has the added benefit of being quicker than culture methods. However, antibiotic sensitivity testing is not possible with LAT alone, so a parallel culture is necessary.
Molecular methods
Polymerase chain reaction (PCR) assays have been proven to be more sensitive than either LAT or culture tests and are highly specific. These PCR tests can be used for capsular typing of encapsulated H. influenzae strains.
Pathogenicity
Host colonization
Many microbes colonize within a host organism. Colonization occurs when a microorganism continues to multiply within the host, without interaction, causing no visible signs of illness or infection. H. influenzae colonizes differently in adults than it does young children. Because this bacterium colonizes more rapidly in young children, they are capable of carrying more than one strain of the same bacterium. Once in the adult stage of life, a human is likely to only be carrying one strain as this bacterium does not colonize as aggressively in adults. Nearly all infants will undergo colonization of this bacteria within their first year of life.
H. influenzae is generally found within and upon the human body, but can also live on various dry, hard surfaces for up to 12 days.
Most strains of H. influenzae are opportunistic pathogens; that is, they usually live in their host without causing disease, but cause problems only when other factors (such as a viral infection, reduced immune function or chronically inflamed tissues, e.g. from allergies) create an opportunity. They infect the host by sticking to the host cell using trimeric autotransporter adhesins.
The pathogenesis of H. influenzae infections is not completely understood, although the presence of the polyribosyl ribitol phosphate (PRP) capsule in encapsulated type b (Hib), a serotype causing conditions such as epiglottitis, is known to be a major factor in virulence. Their capsule allows them to resist phagocytosis and complement-mediated lysis in the nonimmune host. The unencapsulated strains are almost always less invasive; however, they can produce an inflammatory response in humans, which can lead to many symptoms. Vaccination with Hib conjugate vaccine is effective in preventing Hib infection but does not prevent infection with NTHi strains.
H. influenzae can cause respiratory tract infections including pneumonia, otitis media, epiglottitis (swelling in the throat), eye infections and bloodstream infection, meningitis. It can also cause cellulitis (skin infection) and infectious arthritis (inflammation of the joint).
Haemophilus influenzae type b (Hib) infection
Naturally acquired disease caused by H. influenzae seems to occur in humans only. In healthy children under the age of 5, H. influenzae type b was responsible for more than 80% of aggressive infections, before the introduction of the [Hib] vaccine. In infants and young children, H. influenzae type b (Hib) causes bacteremia, pneumonia, epiglottitis and acute bacterial meningitis. On occasion, it causes cellulitis, osteomyelitis, and infectious arthritis. It is one cause of neonatal infection.
Due to routine use of the Hib vaccine in the U.S. since 1990, the incidence of invasive Hib disease has decreased to 1.3/100,000 in children. However, Hib remains a major cause of lower respiratory tract infections in infants and children in developing countries where the vaccine is not widely used. Unencapsulated H. influenzae strains are unaffected by the Hib vaccine and cause ear infections (otitis media), eye infections (conjunctivitis), and sinusitis in children, and are associated with pneumonia.
Treatment
Some strains of H. influenzae produce beta-lactamases, and are also able to modify its penicillin-binding proteins, so the bacteria have gained resistance to the penicillin family of antibiotics. In severe cases, cefotaxime and ceftriaxone delivered directly into the bloodstream are the elected antibiotics, and, for the less severe cases, an association of ampicillin and sulbactam, cephalosporins of the second and third generation, or fluoroquinolones are preferred. (Fluoroquinolone-resistant strains of H. influenzae have been observed).
Macrolides and fluoroquinolones have activity against non-typeable H. influenzae and could be used in patients with a history of allergy to beta-lactam antibiotics. However, macrolide resistance has also been observed.
Serious and chronic complications
The serious complications of HiB are brain damage, hearing loss, and even death. While non-typable H. influenzae strains rarely cause serious disease, they are more likely to cause chronic infections because they have the ability to change their surface antigens. Chronic infections are usually not as serious as acute infections.
There are a few other possible diseases and conditions that can arise from the H. influenzae depending on the areas that they exist in within the human body. This bacterium can exist in the nasal passages (especially the nasopharynx), the ear canal, and the lungs. The bacterium's presence in these areas can lead to some conditions such as otitis media, chronic obstructive pulmonary disorder (COPD), epiglottitis, and asthma which can become severe.
Vaccination
Effective vaccines for Haemophilus influenzae serotype b have been available since the early 1990s, and are recommended for children under age 5 and asplenic patients. The World Health Organization recommends a pentavalent vaccine, combining vaccines against diphtheria, tetanus, pertussis, hepatitis B and Hib. There is not yet sufficient evidence on how effective this pentavalent vaccine is in relation to the individual vaccines.
Hib vaccines cost about seven times the total cost of vaccines against measles, polio, tuberculosis, diphtheria, tetanus, and pertussis. Consequently, whereas 92% of the populations of developed countries were vaccinated against Hib as of 2003, vaccination coverage was 42% for developing countries, and only 8% for least-developed countries.
The Hib vaccines do not provide cross-protection to any other H. influenzae serotypes like Hia, Hic, Hid, Hie or Hif.
An oral vaccination has been developed for non-typeable H. influenzae (NTHi) for patients with chronic bronchitis, but it has not shown to be effective in reducing the number and severity of COPD exacerbations. However, there is no effective vaccine for the other types of capsulated H. influenzae or NTHi.
Vaccines that target unencapsulated H. influenzae serotypes are in development.
| Biology and health sciences | Gram-negative bacteria | Plants |
930016 | https://en.wikipedia.org/wiki/Ascidiacea | Ascidiacea | Ascidiacea, commonly known as the ascidians or sea squirts, is a paraphyletic class in the subphylum Tunicata of sac-like marine invertebrate filter feeders. Ascidians are characterized by a tough outer test or "tunic" made of the polysaccharide cellulose.
Ascidians are found all over the world, usually in shallow water with salinities over 2.5%. While members of the Thaliacea (salps, doliolids and pyrosomes) and Appendicularia (larvaceans) swim freely like plankton, sea squirts are sessile animals after their larval phase: they then remain firmly attached to their substratum, such as rocks and shells.
There are 2,300 species of ascidians and three main types: solitary ascidians, social ascidians that form clumped communities by attaching at their bases, and compound ascidians that consist of many small individuals (each individual is called a zooid) forming large colonies.
Sea squirts feed by taking in water through a tube, the oral siphon. The water enters the mouth and pharynx, flows through mucus-covered gill slits (also called pharyngeal stigmata) into a water chamber called the atrium, then exits through the atrial siphon.
Some authors now include the thaliaceans in Ascidiacea, making it monophyletic.
Anatomy
Sea squirts are rounded or cylindrical animals ranging from about in size. One end of the body is always firmly fixed to rock, coral, or some similar solid surface. The lower surface is pitted or ridged, and in some species has root-like extensions that help the animal grip the surface. The body wall is covered by a smooth thick tunic, which is often quite rigid. The tunic consists of cellulose, along with proteins and calcium salts. Unlike the shells of molluscs, the tunic is composed of living tissue and often has its own blood supply. In some colonial species, the tunics of adjacent individuals are fused into a single structure.
The upper surface of the animal, opposite to the part gripping the substratum, has two openings, or siphons. When removed from the water, the animal often violently expels water from these siphons, hence the common name of "sea squirt". The body itself can be divided into up to three regions, although these are not clearly distinct in most species. The pharyngeal region contains the pharynx, while the abdomen contains most of the other bodily organs, and the postabdomen contains the heart and gonads. In many sea squirts, the postabdomen, or even the entire abdomen, are absent, with their respective organs being located more anteriorly.
As its name implies, the pharyngeal region is occupied mainly by the pharynx. The large buccal siphon opens into the pharynx, acting like a mouth. The pharynx itself is ciliated and contains numerous perforations, or stigmata, arranged in a grid-like pattern around its circumference. The beating of the cilia sucks water through the siphon, and then through the stigmata. A long ciliated groove, or endostyle, runs along one side of the pharynx, and a projecting ridge along the other. The endostyle may be homologous with the thyroid gland of vertebrates, despite its differing function.
The pharynx is surrounded by an atrium, through which water is expelled through a second, usually smaller, siphon. Cords of connective tissue cross the atrium to maintain the general shape of the body. The outer body wall consists of connective tissue, muscle fibres, and a simple epithelium directly underlying the tunic.
Digestive system
The pharynx forms the first part of the digestive system. The endostyle produces a supply of mucus which is then passed into the rest of the pharynx by the beating of flagella along its margins. The mucus then flows in a sheet across the surface of the pharynx, trapping planktonic food particles as they pass through the stigmata, and is collected in the ridge on the dorsal surface. The ridge bears a groove along one side, which passes the collected food downwards and into the oesophageal opening at the base of the pharynx.
The esophagus runs downwards to a stomach in the abdomen, which secretes enzymes that digest the food. An intestine runs upwards from the stomach parallel to the oesophagus and eventually opens, through a short rectum and anus, into a cloaca just below the atrial siphon. In some highly developed colonial species, clusters of individuals may share a single cloaca, with all the atrial siphons opening into it, although the buccal siphons all remain separate. A series of glands lie on the outer surface of the intestine, opening through collecting tubules into the stomach, although their precise function is unclear.
Circulatory system
The heart is a curved muscular tube lying in the postabdomen, or close to the stomach. Each end opens into a single vessel, one running to the endostyle, and the other to the dorsal surface of the pharynx. The vessels are connected by a series of sinuses, through which the blood flows. Additional sinuses run from that on the dorsal surface, supplying blood to the visceral organs, and smaller vessels commonly run from both sides into the tunic. Nitrogenous waste, in the form of ammonia, is excreted directly from the blood through the walls of the pharynx, and expelled through the atrial siphon.
Unusually, the heart of sea squirts alternates the direction in which it pumps blood every three to four minutes. There are two excitatory areas, one at each end of the heart, with first one being dominant, to push the blood through the ventral vessel, and then the other, pushing it dorsally.
There are four different types of blood cell: lymphocytes, phagocytic amoebocytes, nephrocytes and morula cells. The nephrocytes collect waste material such as uric acid and accumulate it in renal vesicles close to the digestive tract. The morula cells help to form the tunic, and can often be found within the tunic substance itself. In some species, the morula cells possess pigmented reducing agents containing iron (hemoglobin), giving the blood a red colour, or vanadium (hemovanadin) giving it a green colour. In that case the cells are also referred to as vanadocytes.
Nervous system
The ascidian central nervous system is formed from a plate that rolls up to form a neural tube. The number of cells within the central nervous system is very small. The neural tube is composed of the sensory vesicle, the neck, the visceral or tail ganglion, and the caudal nerve cord. The anteroposterior regionalization of the neural tube in ascidians is comparable to that in vertebrates.
Although there is no true brain, the largest ganglion is located in the connective tissue between the two siphons, and sends nerves throughout the body. Beneath this ganglion lies an exocrine gland that empties into the pharynx. The gland is formed from the nerve tube, and is therefore homologous to the spinal cord of vertebrates.
Sea squirts lack special sense organs, although the body wall incorporates numerous individual receptors for touch, chemoreception, and the detection of light.
Life history
Almost all ascidians are hermaphrodites and conspicuous mature ascidians are sessile. The gonads are located in the abdomen or postabdomen, and include one testis and one ovary, each of which opens via a duct into the cloaca. Broadly speaking, the ascidians can be divided into species which exist as independent animals (the solitary ascidians) and those which are interdependent (the colonial ascidians). Different species of ascidians can have markedly different reproductive strategies, with colonial forms having mixed modes of reproduction.
Solitary ascidians release many eggs from their atrial siphons; external fertilization in seawater takes place with the coincidental release of sperm from other individuals. A fertilized egg spends 12 hours to a few days developing into a free-swimming tadpole-like larva, which then takes no more than 36 hours to settle and metamorphose into a juvenile.
As a general rule, the larva possesses a long tail, containing muscles, a hollow dorsal nerve tube and a notochord, both features clearly indicative of the animal's chordate affinities. One group though, the molgulid ascidians, have evolved tailless species on at least four separate occasions, and even direct development. A notochord is formed early in development and always consists of a row of exactly 40 cells. The nerve tube enlarges in the main body, and will eventually become the cerebral ganglion of the adult. The tunic develops early in embryonic life and extends to form a fin along the tail in the larva. The larva also has a statocyst and a pigmented cup above the mouth, which opens into a pharynx lined with small clefts opening into a surrounding atrium. The mouth and anus are originally at opposite ends of the animal, with the mouth only moving to its final (posterior) position during metamorphosis.
The larva selects and settles on appropriate surfaces using receptors sensitive to light, orientation to gravity, and tactile stimuli. When its anterior end touches a surface, papillae (small, finger-like nervous projections) secrete an adhesive for attachment. Adhesive secretion prompts an irreversible metamorphosis: various organs (such as the larval tail and fins) are lost while others rearrange to their adult positions, the pharynx enlarges, and organs called ampullae grow from the body to permanently attach the animal to the substratum. The siphons of the juvenile ascidian become orientated to optimise current flow through the feeding apparatus. Sexual maturity can be reached in as little as a few weeks. Since the larva is more advanced than its adult, this type of metamorphosis is called 'retrogressive metamorphosis'. This feature is a landmark for the 'theory of retrogressive metamorphosis or ascidian larva theory'; the true chordates are hypothesized to have evolved from sexually mature larvae.
Direct development in ascidians
Some ascidians, especially in Molgulidae family, have direct development in which the embryo develops directly into the juvenile without developing a tailed larva.
Colonial species
Colonial ascidians reproduce both asexually and sexually. Colonies can survive for decades. An ascidian colony consists of individual elements called zooids. Zooids within a colony are usually genetically identical and some have a shared circulation.
Sexual reproduction
Different colonial ascidian species produce sexually derived offspring by one of two dispersal strategies – colonial species are either broadcast spawners (long-range dispersal) or philopatric (very short-range dispersal). Broadcast spawners release sperm and ova into the water column and fertilization occurs near to the parent colonies. Some species are also viviparous. Resultant zygotes develop into microscopic larvae that may be carried great distances by oceanic currents. The larvae of sessile forms which survive eventually settle and complete maturation on the substratum- then they may bud asexually to form a colony of zooids.
The picture is more complicated for the philopatrically dispersed ascidians: sperm from a nearby colony (or from a zooid of the same colony) enter the atrial siphon and fertilization takes place within the atrium. Embryos are then brooded within the atrium where embryonic development takes place: this results in macroscopic tadpole-like larvae. When mature, these larvae exit the atrial siphon of the adult and then settle close to the parent colony (often within meters). The combined effect of short sperm range and philopatric larval dispersal results in local population structures of closely related individuals/inbred colonies. Generations of colonies which are restricted in dispersal are thought to accumulate adaptions to local conditions, thereby providing advantages over newcomers.
Trauma or predation often results in fragmentation of a colony into subcolonies. Subsequent zooid replication can lead to coalescence and circulatory fusion of the subcolonies. Closely related colonies which are proximate to each other may also fuse if they coalesce and if they are histocompatible. Ascidians were among the first animals to be able to immunologically distinguish self from non-self as a mechanism to prevent unrelated colonies from fusing to them and parasitizing them.
Fertilization
Sea squirt eggs are surrounded by a fibrous vitelline coat and a layer of follicle cells that produce sperm-attracting substances. In fertilization, the sperm passes through the follicle cells and binds to glycosides on the vitelline coat. The sperm's mitochondria are left behind as the sperm enters and drives through the coat; this translocation of the mitochondria might provide the necessary force for penetration. The sperm swims through the perivitelline space, finally reaching the egg plasma membrane and entering the egg. This prompts rapid modification of the vitelline coat, through processes such as the egg's release of glycosidase into the seawater, so no more sperm can bind and polyspermy is avoided. After fertilization, free calcium ions are released in the egg cytoplasm in waves, mostly from internal stores. The temporary large increase in calcium concentration prompts the physiological and structural changes of development.
The dramatic rearrangement of egg cytoplasm following fertilization, called ooplasmic segregation, determines the dorsoventral and anteroposterior axes of the embryo. There are at least three types of sea squirt egg cytoplasm: ectoplasm containing vesicles and fine particles, endoderm containing yolk platelets, and myoplasm containing pigment granules, mitochondria, and endoplasmic reticulum. In the first phase of ooplasmic segregation, the myoplasmic actin-filament network contracts to rapidly move the peripheral cytoplasm (including the myoplasm) to the vegetal pole, which marks the dorsal side of the embryo. In the second phase, the myoplasm moves to the subequatorial zone and extends into a crescent, which marks the future posterior of the embryo. The ectoplasm with the zygote nucleus ends up at the animal hemisphere while the endoplasm ends up in the vegetal hemisphere.
Promotion of out-crossing
Ciona intestinalis is a hermaphrodite that releases sperm and eggs into the surrounding seawater almost simultaneously. It is self-sterile, and thus has been used for studies on the mechanism of self-incompatibility. Self/non-self-recognition molecules play a key role in the process of interaction between sperm and the vitelline coat of the egg. It appears that self/non-self recognition in ascidians such as C. intestinalis is mechanistically similar to self-incompatibility systems in flowering plants. Self-incompatibility promotes out-crossing, and thus provides the adaptive advantage at each generation of masking deleterious recessive mutations (i.e. genetic complementation).
Ciona savignyi is highly self-fertile. However, non-self sperm out-compete self-sperm in fertilization competition assays. Gamete recognition is not absolute allowing some self-fertilization. It was speculated that self-incompatibility evolved to avoid inbreeding depression, but that selfing ability was retained to allow reproduction at low population density.
Botryllus schlosseri is a colonial tunicate able to reproduce both sexually and asexually. B. schlosseri is a sequential (protogynous) hermaphrodite, and in a colony, eggs are ovulated about two days before the peak of sperm emission. Thus self-fertilization is avoided, and cross-fertilization is favored. Although avoided, self-fertilization is still possible in B. schlosseri. Self-fertilized eggs develop with a substantially higher frequency of anomalies during cleavage than cross-fertilized eggs (23% vs. 1.6%). Also, a significantly lower percentage of larvae derived from self-fertilized eggs metamorphose, and the growth of the colonies derived from their metamorphosis is significantly lower. These findings suggest that self-fertilization gives rise to inbreeding depression associated with developmental deficits that are likely caused by expression of deleterious recessive mutations.
Asexual reproduction
Many colonial sea squirts are also capable of asexual reproduction, although the means of doing so are highly variable between different families. In the simplest forms, the members of the colony are linked only by rootlike projections from their undersides known as stolons. Buds containing food storage cells can develop within the stolons and, when sufficiently separated from the 'parent', may grow into a new adult individual.
In other species, the postabdomen can elongate and break up into a string of separate buds, which can eventually form a new colony. In some, the pharyngeal part of the animal degenerates, and the abdomen breaks up into patches of germinal tissue, each combining parts of the epidermis, peritoneum, and digestive tract, and capable of growing into new individuals.
In yet others, budding begins shortly after the larva has settled onto the substrate. In the family Didemnidae, for instance, the individual essentially splits into two, with the pharynx growing a new digestive tract and the original digestive tract growing a new pharynx.
DNA repair
Apurinic/apyrimidinic (AP) sites are a common form of DNA damage that inhibit DNA replication and transcription. AP endonuclease 1 (APEX1), an enzyme produced by C. intestinalis, is employed in the repair of AP sites during early embryonic development. Lack of such repair leads to abnormal development. C. intestinalis also has a set of genes that encode proteins homologous to those employed in the repair of DNA interstrand crosslinks in humans.
Ecology
The exceptional filtering capability of adult sea squirts causes them to accumulate pollutants that may be toxic to embryos and larvae as well as impede enzyme function in adult tissues. This property has made some species sensitive indicators of pollution.
Over the last few hundred years, most of the world's harbors have been invaded by non-native sea squirts that have been introduced by accident from the shipping industry. Several factors, including quick attainment of sexual maturity, tolerance of a wide range of environments, and a lack of predators, allow sea squirt populations to grow rapidly. Unwanted populations on docks, ship hulls, and farmed shellfish cause significant economic problems, and sea squirt invasions have disrupted the ecosystem of several natural sub-tidal areas by smothering native animal species.
Sea squirts are the natural prey of many animals, including nudibranchs, flatworms, molluscs, rock crabs, sea stars, fish, birds, and sea otters. Some are also eaten by humans in many parts of the world, including Japan, Korea, Chile, and Europe (where they are sold under the name "sea violet"). As chemical defenses, many sea squirts intake and maintain an extremely high concentration of vanadium in the blood, have a very low pH of the tunic due to acid in easily ruptured bladder cells, and (or) produce secondary metabolites harmful to predators and invaders. Some of these metabolites are toxic to cells and are of potential use in pharmaceuticals.
Evolution
Fossil record
Ascidians are soft-bodied animals, and for this reason, their fossil record is almost entirely lacking. The earliest reliable ascidians is Shankouclava shankouense from the Lower Cambrian Maotianshan Shale (Yunnan, South China). There are also two enigmatic species from the Ediacaran period with some affinity to the ascidians – Ausia from the Nama Group of Namibia and Burykhia from the Onega Peninsula, White Sea of northern Russia.
They are also recorded from Lower Jurassic (Bonet and Benveniste-Velasquez, 1971; Buge and Monniot, 1972) and the Tertiary from France (Deflandre-Riguard, 1949, 1956; Durand, 1952; Deflandre and Deflandre-Rigaud, 1956; Bouche, 1962; Lezaud, 1966; Monniot and Buge, 1971; Varol and Houghton, 1996). Older (Triassic) records are ambiguous. From the Early Jurassic, the species Didemnum cassianum, Quadrifolium hesselboi, Palaeoquadrum ullmanni and other indet genera are recorded. The representatives of the genus Cystodytes (family Polycitoridae) have been described from the Pliocene of France by Monniot (1970, 1971) and Deflandre-Rigaud (1956), and from Eocene of France by Monniot and Buge (1971), and lately from the Late Eocene of S Australia by Łukowiak (2012).
Phylogeny
The ascidians were on morphological evidence treated as sister to the Thaliacea and Appendicularia, but molecular evidence has suggested that ascidians could be polyphyletic within the Tunicata, as shown in the following cladogram.
In 2017 and 2018, two studies were published, which suggested an alternate phylogeny, placing Appendicularia as sister to the rest of Tunicata, and Thaliacea nested inside Ascidiacea. A grouping of Thaliacea and Ascidiacea to the exclusion of Appendicularia had already been suggested for a long time, under the name of Acopa. Brusca et al. treat Ascidiacea as a monophyletic group including pelagic Thaliacea.
Uses
Culinary
Various ascidians are eaten by humans around the world as delicacies.
Sea pineapple (Halocynthia roretzi) is cultivated in Japan (hoya, maboya) and Korea (meongge). When served raw, they have a chewy texture and peculiar flavor likened to "rubber dipped in ammonia" which has been attributed to a naturally occurring chemical known as cynthiaol. Styela clava is farmed in parts of Korea where it is known as mideoduk and is added to various seafood dishes such as agujjim. Tunicate bibimbap is a specialty of Geoje Island, not far from Masan.
Microcosmus species from the Mediterranean Sea are eaten in France (figue de mer, violet), Italy (limone di mare, uova di mare) and Greece (fouska, φούσκα), for example, raw with lemon, or in salads with olive oil, lemon and parsley.
The piure (Pyura chilensis) is used in the cuisine of Chile – it is consumed both raw and in seafood stews similar to bouillabaisse.
Pyura praeputialis is known as cunjevoi in Australia. It was once used as a food source by Aboriginal people living around Botany Bay, but is now used mainly for fishing bait.
Ciona is being developed in Norway as a potential substitute meat protein, after processing to remove its 'marine taste' and to make its texture less 'squid-like'.
Model organisms for research
Several factors make sea squirts good models for studying the fundamental developmental processes of chordates, such as cell-fate specification. The embryonic development of sea squirts is simple, rapid, and easily manipulated. Because each embryo contains relatively few cells, complex processes can be studied at the cellular level, while remaining in the context of the whole embryo.
The eggs of some species contain little yolk and are therefore transparent making them transparency ideal for fluorescent imaging. Its maternally-derived proteins are naturally associated with pigment (in a few species only), so cell lineages are easily labeled, allowing scientists to visualize embryogenesis from beginning to end.
Sea squirts are also valuable because of their unique evolutionary position: as an approximation of ancestral chordates, they can provide insight into the link between chordates and ancestral non-chordate deuterostomes, as well as the evolution of vertebrates from simple chordates. The sequenced genomes of the related sea squirts Ciona intestinalis and Ciona savignyi are small and easily manipulated; comparisons with the genomes of other organisms such as flies, nematodes, pufferfish and mammals provides valuable information regarding chordate evolution. A collection of over 480,000 cDNAs have been sequenced and are available to support further analysis of gene expression, which is expected to provide information about complex developmental processes and regulation of genes in vertebrates. Gene expression in embryos of sea squirts can be conveniently inhibited using Morpholino oligos.
| Biology and health sciences | Chordata (except vertebrates) | Animals |
25853144 | https://en.wikipedia.org/wiki/Frecciarossa | Frecciarossa | Frecciarossa (; from , "red arrow") is a high-speed train of the Italian national train operator, Trenitalia, as well as a member of the train category Le Frecce. The name was introduced in 2008 after it had previously been known as Eurostar Italia. Frecciarossa trains operate at speeds of up to . Frecciarossa is the premier service of Trenitalia and competes with italo, operated by Nuovo Trasporto Viaggiatori. Trenitalia also operates the sister brands Frecciargento and Frecciabianca for slower services.
Routes
Frecciarossa trains travel on dedicated high-speed railway lines and, on some routes, also on conventional railway lines with lower speed limits. Current limitations on the tracks set the maximum operating speed of both types of trains to . Frecciarossa trains operate the following services:
Turin - Milan - Reggio Emilia AV - Bologna - Florence - Rome - Naples - Salerno
Turin - Milan - Brescia - Verona - Vicenza - Padua - Venice - Monfalcone - Trieste
Venice - Padua - Bologna - Florence - Rome - Naples - Salerno
Bergamo - Brescia - Verona - Bologna - Florence - Rome
Udine - Pordenone - Treviso - Venice - Padua - Vicenza - Verona - Brescia - Milan
Milan - Reggio Emilia AV - Bologna - Rimini - Ancona - S. Benedetto T. - Pescara - Termoli - Foggia - Bari - Brindisi - Lecce
Milan - Bologna - Florence - Rome - Naples - Salerno - Potenza - Ferrandina - Metaponto - Taranto
Venice - Padua - Vicenza - Verona - Brescia - Milan - Pavia - Genoa
Venice - Padua - Ferrara - Bologna - Florence - Rome - Naples - Salerno
Perugia - Arezzo - Florence - Bologna - Reggio Emilia AV - Milan - Turin
Milan - Reggio Emilia AV - Bologna - Florence - Rome - Naples - Salerno - Agropoli - Sapri
Milan - Reggio Emilia EV - Bologna - Florence - Paola - Lamezia - Rosarno - Villa San Giovanni - Reggio Calabria
The brand also includes the Milan–Paris Frecciarossa, which operates two routes:
Milan – Turin – Bardonecchia (seasonal) – Modane – Chambéry-Challes-les-Eaux – Lyon-Part-Dieu – Paris Gare de Lyon
Lyon-Perrache – Lyon-Part-Dieu – Paris Gare de Lyon
Rolling stock
The following rolling stock types are used for Frecciarossa services:
ETR.500: non-tilting train made of eleven passenger coaches (one with cafe/restaurant service) with 574 seats moved by two E.404 locomotives, speeds up to .
ETR.600: tilting train made of seven passenger coaches (one with cafe/restaurant service) with 432 seats, speeds up to 250 km/h (155 mph).
ETR.700: non-tilting train made of 8 passenger coaches (one with cafe/restaurant service) with 497 seats, speeds up to 250 km/h (155 mph).
ETR.1000: non-tilting electro-train made of eight passenger coaches (one with cafe/restaurant service) with 457 seats, speeds up to .
Accidents and incidents
On 6 February 2020, a Frecciarossa train derailed at Ospedaletto Lodigiano, killing two people and injuring 27 others.
On 10 December 2023, a Frecciarossa train collided with another passenger train at Faenza injuring 17.
| Technology | High-speed rail | null |
2875647 | https://en.wikipedia.org/wiki/Energy%E2%80%93momentum%20relation | Energy–momentum relation | In physics, the energy–momentum relation, or relativistic dispersion relation, is the relativistic equation relating total energy (which is also called relativistic energy) to invariant mass (which is also called rest mass) and momentum. It is the extension of mass–energy equivalence for bodies or systems with non-zero momentum.
It can be formulated as:
This equation holds for a body or system, such as one or more particles, with total energy , invariant mass , and momentum of magnitude ; the constant is the speed of light. It assumes the special relativity case of flat spacetime and that the particles are free. Total energy is the sum of rest energy and relativistic kinetic energy:
Invariant mass is mass measured in a center-of-momentum frame.
For bodies or systems with zero momentum, it simplifies to the mass–energy equation , where total energy in this case is equal to rest energy.
The Dirac sea model, which was used to predict the existence of antimatter, is closely related to the energy–momentum relation.
Connection to E = mc2
The energy–momentum relation is consistent with the familiar mass–energy relation in both its interpretations: relates total energy to the (total) relativistic mass (alternatively denoted or ), while relates rest energy to (invariant) rest mass .
Unlike either of those equations, the energy–momentum equation () relates the total energy to the rest mass . All three equations hold true simultaneously.
Special cases
If the body is a massless particle (), then () reduces to . For photons, this is the relation, discovered in 19th century classical electromagnetism, between radiant momentum (causing radiation pressure) and radiant energy.
If the body's speed is much less than , then () reduces to ; that is, the body's total energy is simply its classical kinetic energy () plus its rest energy.
If the body is at rest (), i.e. in its center-of-momentum frame (), we have and ; thus the energy–momentum relation and both forms of the mass–energy relation (mentioned above) all become the same.
A more general form of relation () holds for general relativity.
The invariant mass (or rest mass) is an invariant for all frames of reference (hence the name), not just in inertial frames in flat spacetime, but also accelerated frames traveling through curved spacetime (see below). However the total energy of the particle and its relativistic momentum are frame-dependent; relative motion between two frames causes the observers in those frames to measure different values of the particle's energy and momentum; one frame measures and , while the other frame measures and , where and , unless there is no relative motion between observers, in which case each observer measures the same energy and momenta. Although we still have, in flat spacetime:
The quantities , , , are all related by a Lorentz transformation. The relation allows one to sidestep Lorentz transformations when determining only the magnitudes of the energy and momenta by equating the relations in the different frames. Again in flat spacetime, this translates to;
Since does not change from frame to frame, the energy–momentum relation is used in relativistic mechanics and particle physics calculations, as energy and momentum are given in a particle's rest frame (that is, and as an observer moving with the particle would conclude to be) and measured in the lab frame (i.e. and as determined by particle physicists in a lab, and not moving with the particles).
In relativistic quantum mechanics, it is the basis for constructing relativistic wave equations, since if the relativistic wave equation describing the particle is consistent with this equation – it is consistent with relativistic mechanics, and is Lorentz invariant. In relativistic quantum field theory, it is applicable to all particles and fields.
Origins and derivation of the equation
The energy–momentum relation goes back to Max Planck's article
published in 1906.
It was used by Walter Gordon in 1926 and then by Paul Dirac in 1928 under the form , where V is the amount of potential energy.
The equation can be derived in a number of ways, two of the simplest include:
From the relativistic dynamics of a massive particle,
By evaluating the norm of the four-momentum of the system. This method applies to both massive and massless particles, and can be extended to multi-particle systems with relatively little effort (see below).
Heuristic approach for massive particles
For a massive object moving at three-velocity with magnitude in the lab frame:
is the total energy of the moving object in the lab frame,
is the three dimensional relativistic momentum of the object in the lab frame with magnitude . The relativistic energy and momentum include the Lorentz factor defined by:
Some authors use relativistic mass defined by:
although rest mass has a more fundamental significance, and will be used primarily over relativistic mass in this article.
Squaring the 3-momentum gives:
then solving for and substituting into the Lorentz factor one obtains its alternative form in terms of 3-momentum and mass, rather than 3-velocity:
Inserting this form of the Lorentz factor into the energy equation gives:
followed by more rearrangement it yields (). The elimination of the Lorentz factor also eliminates implicit velocity dependence of the particle in (), as well as any inferences to the "relativistic mass" of a massive particle. This approach is not general as massless particles are not considered. Naively setting would mean that and and no energy–momentum relation could be derived, which is not correct.
Norm of the four-momentum
Special relativity
In Minkowski space, energy (divided by ) and momentum are two components of a Minkowski four-vector, namely the four-momentum;
(these are the contravariant components).
The Minkowski inner product of this vector with itself gives the square of the norm of this vector, it is proportional to the square of the rest mass of the body:
a Lorentz invariant quantity, and therefore independent of the frame of reference. Using the Minkowski metric with metric signature , the inner product is
and
so
or, in natural units where = 1,
General relativity
In general relativity, the 4-momentum is a four-vector defined in a local coordinate frame, although by definition the inner product is similar to that of special relativity,
in which the Minkowski metric is replaced by the metric tensor field :
solved from the Einstein field equations. Then:
Units of energy, mass and momentum
In natural units where , the energy–momentum equation reduces to
In particle physics, energy is typically given in units of electron volts (eV), momentum in units of eV·−1, and mass in units of eV·−2. In electromagnetism, and because of relativistic invariance, it is useful to have the electric field and the magnetic field in the same unit (Gauss), using the cgs (Gaussian) system of units, where energy is given in units of erg, mass in grams (g), and momentum in g·cm·s−1.
Energy may also in theory be expressed in units of grams, though in practice it requires a large amount of energy to be equivalent to masses in this range. For example, the first atomic bomb liberated about 1 gram of heat, and the largest thermonuclear bombs have generated a kilogram or more of heat. Energies of thermonuclear bombs are usually given in tens of kilotons and megatons referring to the energy liberated by exploding that amount of trinitrotoluene (TNT).
Special cases
Centre-of-momentum frame (one particle)
For a body in its rest frame, the momentum is zero, so the equation simplifies to
where is the rest mass of the body.
Massless particles
If the object is massless, as is the case for a photon, then the equation reduces to
This is a useful simplification. It can be rewritten in other ways using the de Broglie relations:
if the wavelength or wavenumber are given.
Correspondence principle
Rewriting the relation for massive particles as:
and expanding into power series by the binomial theorem (or a Taylor series):
in the limit that , we have so the momentum has the classical form , then to first order in (i.e. retain the term for and neglect all terms for ) we have
or
where the second term is the classical kinetic energy, and the first is the rest energy of the particle. This approximation is not valid for massless particles, since the expansion required the division of momentum by mass. Incidentally, there are no massless particles in classical mechanics.
Many-particle systems
Addition of four momenta
In the case of many particles with relativistic momenta and energy , where (up to the total number of particles) simply labels the particles, as measured in a particular frame, the four-momenta in this frame can be added;
and then take the norm; to obtain the relation for a many particle system:
where is the invariant mass of the whole system, and is not equal to the sum of the rest masses of the particles unless all particles are at rest (see for more detail). Substituting and rearranging gives the generalization of ();
The energies and momenta in the equation are all frame-dependent, while is frame-independent.
Center-of-momentum frame
In the center-of-momentum frame (COM frame), by definition we have:
with the implication from () that the invariant mass is also the centre of momentum (COM) mass–energy, aside from the factor:
and this is true for all frames since is frame-independent. The energies are those in the COM frame, not the lab frame. However, many familiar bound systems have the lab frame as COM frame, since the system itself is not in motion and so the momenta all cancel to zero. An example would be a simple object (where vibrational momenta of atoms cancel) or a container of gas where the container is at rest. In such systems, all the energies of the system are measured as mass. For example, the heat in an object on a scale, or the total of kinetic energies in a container of gas on the scale, all are measured by the scale as the mass of the system.
Rest masses and the invariant mass
Either the energies or momenta of the particles, as measured in some frame, can be eliminated using the energy momentum relation for each particle:
allowing to be expressed in terms of the energies and rest masses, or momenta and rest masses. In a particular frame, the squares of sums can be rewritten as sums of squares (and products):
so substituting the sums, we can introduce their rest masses in ():
The energies can be eliminated by:
similarly the momenta can be eliminated by:
where is the angle between the momentum vectors and .
Rearranging:
Since the invariant mass of the system and the rest masses of each particle are frame-independent, the right hand side is also an invariant (even though the energies and momenta are all measured in a particular frame).
Matter waves
Using the de Broglie relations for energy and momentum for matter waves,
where is the angular frequency and is the wavevector with magnitude , equal to the wave number, the energy–momentum relation can be expressed in terms of wave quantities:
and tidying up by dividing by throughout:
This can also be derived from the magnitude of the four-wavevector
in a similar way to the four-momentum above.
Since the reduced Planck constant and the speed of light both appear and clutter this equation, this is where natural units are especially helpful. Normalizing them so that , we have:
Tachyon and exotic matter
The velocity of a bradyon with the relativistic energy–momentum relation
can never exceed . On the contrary, it is always greater than for a tachyon whose energy–momentum equation is
By contrast, the hypothetical exotic matter has a negative mass and the energy–momentum equation is
| Physical sciences | Theory of relativity | Physics |
2876360 | https://en.wikipedia.org/wiki/Potassium%20manganate | Potassium manganate | Potassium manganate is the inorganic compound with the formula . This green-colored salt is an intermediate in the industrial synthesis of potassium permanganate (), a common chemical. Occasionally, potassium manganate and potassium permanganate are confused, but each compound's properties are distinct.
Structure and bonding
is a salt, consisting of cations and anions. X-ray crystallography shows that the anion is tetrahedral, with Mn-O distances of 1.66 Å, ca. 0.03 Å longer than the Mn-O distances in . It is isostructural with potassium sulfate. The compound is paramagnetic, owing to the presence of one unpaired electron on the Mn(VI) center.
Synthesis
The industrial route entails treatment of with air and potassium hydroxide:
2 MnO2 + 4 KOH + O2 → 2 K2MnO4 + 2 H2O
The transformation gives a green-colored melt. Alternatively, instead of using air, potassium nitrate can be used as the oxidizer:
One can test an unknown substance for the presence of manganese by heating the sample in strong KOH in air. The production of a green coloration indicates the presence of Mn. This green color results from an intense absorption at 610 nm.
In the laboratory, can be synthesized by heating a solution of in concentrated KOH solution followed by cooling to give green crystals:
This reaction illustrates the relatively rare role of hydroxide as a reducing agent. The concentration of in such solutions can be checked by measuring their absorbance at 610 nm.
The one-electron reduction of permanganate to manganate can also be effected using iodide as the reducing agent:
The conversion is signaled by the color change from purple, characteristic of permanganate, to the green color of manganate. This reaction also shows that manganate(VII) can serve as an electron acceptor in addition to its usual role as an oxygen-transfer reagent. Barium manganate, , is generated by the reduction of with iodide in the presence of barium chloride. Just like , exhibits low solubility in virtually all solvents.
An easy method for preparing potassium manganate in the laboratory involves heating crystals or powder of pure potassium permanganate. Potassium permanganate will decompose into potassium manganate, manganese dioxide and oxygen gas:
This reaction is a laboratory method to prepare oxygen, but produces samples of potassium manganate contaminated with . The former is soluble and the latter is not.
Reactions
Manganate salts readily disproportionate to permanganate ion and manganese dioxide:
The colorful nature of the disproportionation has led the manganate/manganate(VII) pair to be referred to as a chemical chameleon. This disproportionation reaction, which becomes rapid when [] < 1M, follows bimolecular kinetics.
| Physical sciences | Metallic oxyanions | Chemistry |
38540045 | https://en.wikipedia.org/wiki/Aqueduct%20%28water%20supply%29 | Aqueduct (water supply) | An aqueduct is a watercourse constructed to carry water from a source to a distribution point far away. In modern engineering, the term aqueduct is used for any system of pipes, ditches, canals, tunnels, and other structures used for this purpose. The term aqueduct also often refers specifically to a bridge carrying an artificial watercourse.
Aqueducts were used in ancient Greece, the ancient Near East, ancient Rome, ancient Aztec, and ancient Inca. The simplest aqueducts are small ditches cut into the earth. Much larger channels may be used in modern aqueducts. Aqueducts sometimes run for some or all of their path through tunnels constructed underground. Modern aqueducts may also use pipelines. Historically, agricultural societies have constructed aqueducts to irrigate crops and supply large cities with drinking water.
Etymology
The word aqueduct is derived from the Latin words (water) and (led or guided).
Ancient aqueducts
Although particularly associated with the Romans, aqueducts were devised much earlier in Greece, the Near East, Nile Valley, and Indian subcontinent, where people such as the Egyptians and Harappans built sophisticated irrigation systems. The Aztecs and Incans also built such systems independently later. Roman-style aqueducts were used as early as the 7th century BC, when the Assyrians built an 80 km long limestone aqueduct, which included a 10 m high section to cross a 300 m wide valley, to carry water to their capital city, Nineveh.
Crete
Although particularly associated with the Romans, aqueducts were likely first used by the Minoans around 2000 BC. The Minoans had developed what was then an extremely advanced irrigation system, including several aqueducts.
Deccan
The Indian subcontinent is believed to have some of the earliest aqueducts. Evidence can be found at the sites of present-day Hampi, Karnataka. The massive aqueducts near Tungabhadra River supplying irrigation water were once long. The waterways supplied water to royal bath tubs.
Oman
In Oman from the Iron Age, in Salut, Bat, and other sites, a system of underground aqueducts called falaj or qanāts were constructed, a series of well-like vertical shafts, connected by gently sloping horizontal tunnels.
There are three types of falaj:
Daudi (داوودية) with underground aqueducts
Ghaili (الغيلية ) requiring a dam to collect the water
Aini (العينية ) whose source is a water spring
These enabled large scale agriculture to flourish in a dry land environment.
Persia
In Persia, starting around 3000 years ago a system of underground aqueducts called qanāts were constructed, a series of well-like vertical shafts, connected by gently sloping tunnels. This technique:
taps into subterranean water in a manner that delivers water to the surface without the need for pumping. The water drains relied on gravity, with the destination lower than the source, which is typically an upland aquifer.
allows water to be transported long distances in hot dry climates without losing a large proportion of the source water to seepage and evaporation.
Petra, Jordan
Throughout Petra, Jordan, the Nabataean engineers took advantage of every natural spring and every winter downpour to channel water where it was needed. They constructed aqueducts and piping systems that allowed water to flow across mountains, through gorges and into the temples, homes, and gardens of Petra's citizens. Walking through the Siq, one can easily spot the remains of channels that directed water to the city center, as well as durable retention dams that kept powerful flood waters at bay.
Greece
On the island of Samos, the Tunnel of Eupalinos was built during the reign of Polycrates (538–522 BC). It is considered an underground aqueduct and brought fresh water to Pythagoreion for roughly a thousand years.
Roman
Roman aqueducts were built in all parts of the Roman Empire, from Germany to Africa, and especially in the city of Rome, where they totalled over . The aqueducts supplied fresh water to public baths and for drinking water, in large cities across the empire, and set a standard of engineering that was not surpassed for more than a thousand years. Bridges, built in stone with multiple arches, were a distinctive feature of Roman aqueducts and hence the term aqueduct is often applied specifically to a bridge for carrying water.
South America
Near the Peruvian town of Nazca, an ancient pre-Columbian system of aqueducts called puquios were built and are still in use today. They were made of intricately placed stones, a construction material widely used by the Nazca culture. The time period in which they were constructed is still debated, but some evidence supports circa A.D. 540–552, in response to drought periods in the region.
The Guayabo National Monument of Costa Rica, a park covering the largest archaeological site in the country, contains a system of aqueducts. The complex network of uncovered and covered aqueducts still functions well. The aqueducts are constructed from rounded river stones, which are mostly made of volcanic rock. The civilization that constructed the aqueduct system remains a mystery to archaeologists; it is suspected that Guayabo's aqueducts sat at a point of ancient cultural confluence between Aztecs, Mayans, and Incas.
North America
When Europeans saw the Aztec capital Tenochtitlan, early in the 16th century, the city was watered by two aqueducts. One of these, Chapultepec aqueduct, built , was rebuilt by the Spanish almost three hundred years later. Originally tracing part of its path over now-gone Lake Texcoco, only a fragment remains in Mexico City today.
Sri Lanka
Extensive usage of elaborate aqueducts have been found to have been used in ancient Sri Lanka. The best example is the Yoda Ela or Jaya Ganga, an long water canal carrying excess water between two artificial reservoirs with a gradient of 10 to 20 cm per kilometer during the fifth century AD. However, the ancient engineering methods in calculating the exact elevation between the two reservoirs and the exact gradient of the canal to such fine precision had been lost with the fall of the civilization in 13th Century.
Modern aqueducts
Modern aqueducts are a central part of many countries' water distribution infrastructure.
The United States' aqueducts are some of the world's largest. The Catskill Aqueduct carries water to New York City over a distance of 120 miles (190 km), but is dwarfed by aqueducts in the far west of the country, most notably the 242-mile (389-km) Colorado River Aqueduct, which supplies the Los Angeles area with water from the Colorado River nearly 250 miles to the east and the California Aqueduct, which runs from the Sacramento-San Joaquin River Delta to Lake Perris. The Central Arizona Project is the largest and most expensive aqueduct constructed in the United States. It stretches 336 miles from its source near Parker, Arizona to the metropolitan areas of Phoenix and Tucson.
An aqueduct in New Zealand, "the Oamaru Borough Race", was constructed in the late 19th century to deliver water (and water-power) about 50 km from the Waitaki River at Kurow to the coastal town of Oamaru.
In Spain, the Tagus-Segura Water Transfer system of aqueducts opened in 1979 and transports water from north to south.
In China, the South–North Water Transfer Project aims to connect the Yangtze River basin to Beijing through three separate systems. The project will reuse part of the Grand Canal of China.
Design
Open channels
The simplest aqueducts are small ditches cut into the earth. Much larger channels may be used in modern aqueducts, for instance the Central Arizona Project uses wide channels. A major factor in the design of all open channels is its gradient. A higher gradient allows a smaller channel to carry the same amount of water as a larger channel with a lower gradient, but increases the potential of the water to damage the aqueduct's structure. A typical Roman aqueduct had a gradient of about 1:4800.
Artificial rills
A constructed functional rill is a small canal or aqueduct of stone, brick, concrete, or other lining material, usually rectilinear in cross section, for water transportation from a source such as a river, spring, reservoir, qanat, or aqueduct for domestic consumption or agricultural irrigation of crop land uses.
Rills were traditionally used in Middle Eastern and Mediterranean climate cultures of ancient and historical eras; and other climates and continents worldwide. They are distinguished from a 'water ditch' by being lined to reduce absorption losses and to increase durability. The Falaj irrigation system at the Al Ain Oasis, in present-day Abu Dhabi Emirate, uses rills as part of its qanat water system. Sometimes in the Spanish language they are called Acequias.
Rills are also used for aesthetic purposes in landscape design. Rills are used as narrow channels of water inset into the pavement of a garden, as linear water features, and often tiled and part of a fountain design.
The historical origins are from paradise garden religious images that first translated into ancient Persian Gardens. Rills were later exceptionally developed in the Moorish (Spanish) Gardens of Al-andalus, such as at the Alhambra in Granada; and also in other Islamic gardens, cultures, and countries. Early 20th century examples are in the María Luisa Park gardens in Seville, Spain; and at the Casa del Herrero gardens in Montecito, California.
Tunnels
Aqueducts sometimes run for some or all of their path through tunnels constructed underground. A version of this common in North Africa and Central Asia that has vertical wells at regular intervals is called a qanat. One historic example found in Syria, the Qanat Firaun, extends over 100 kilometers.
Pipes
Modern aqueducts may also make extensive use of pipelines. Pipelines are useful for transporting water over long distances when it needs to move over hills, or where open channels are poor choices due to considerations of evaporation, freezing, pollution, or environmental impact. They can also be used to carry treated water.
Uses
Historically, agricultural societies have constructed aqueducts to irrigate crops. Archimedes invented the water screw to raise water for use in irrigation of croplands.
Another use for aqueducts is to supply large cities with drinking water. They also help drought-prone areas with water supply. Some of the Roman aqueducts still supply water to Rome today. In California, United States, three large aqueducts supply water over hundreds of miles to the Los Angeles area. Two are from the Owens River area, and a third is from the Colorado River.
In modern civil engineering projects, detailed study and analysis of open-channel flow is commonly required to support flood control, irrigation systems, and large water supply systems when an aqueduct rather than a pipeline is the preferred solution.
In the past, aqueducts often had channels made of earth or other porous materials but significant amounts of water are lost through such unlined aqueducts. As water gets increasingly scarce, these canals are being lined with concrete, polymers, or impermeable soil. In some cases, a new aqueduct is built alongside the old one because it cannot be shut down during construction.
| Technology | Food, water and health | null |
28996847 | https://en.wikipedia.org/wiki/Geophysical%20fluid%20dynamics | Geophysical fluid dynamics | Geophysical fluid dynamics, in its broadest meaning, refers to the fluid dynamics of naturally occurring flows, such as lava flows, oceans, and planetary atmospheres, on Earth and other planets.
Two physical features that are common to many of the phenomena studied in geophysical fluid dynamics are rotation of the fluid due to the planetary rotation and stratification (layering). The applications of geophysical fluid dynamics do not generally include the circulation of the mantle, which is the subject of geodynamics, or fluid phenomena in the magnetosphere.
Fundamentals
To describe the flow of geophysical fluids, equations are needed for conservation of momentum (or Newton's second law) and conservation of energy. The former leads to the Navier–Stokes equations which cannot be solved analytically (yet). Therefore, further approximations are generally made in order to be able to solve these equations. First, the fluid is assumed to be incompressible. Remarkably, this works well even for a highly compressible fluid like air as long as sound and shock waves can be ignored. Second, the fluid is assumed to be a Newtonian fluid, meaning that there is a linear relation between the shear stress and the strain , for example
where is the viscosity. Under these assumptions the Navier-Stokes equations are
The left hand side represents the acceleration that a small parcel of fluid would experience in a reference frame that moved with the parcel (a Lagrangian frame of reference). In a stationary (Eulerian) frame of reference, this acceleration is divided into the local rate of change of velocity and advection, a measure of the rate of flow in or out of a small region.
The equation for energy conservation is essentially an equation for heat flow. If heat is transported by conduction, the heat flow is governed by a diffusion equation. If there are also buoyancy effects, for example hot air rising, then natural convection, also known as free convection, can occur. Convection in the Earth's outer core drives the geodynamo that is the source of the Earth's magnetic field. In the ocean, convection can be thermal (driven by heat), haline (where the buoyancy is due to differences in salinity), or thermohaline, a combination of the two.
Buoyancy and stratification
Fluid that is less dense than its surroundings tends to rise until it has the same density as its surroundings. If there is not much energy input to the system, it will tend to become stratified. On a large scale, Earth's atmosphere is divided into a series of layers. Going upwards from the ground, these are the troposphere, stratosphere, mesosphere, thermosphere, and exosphere.
The density of air is mainly determined by temperature and water vapor content, the density of sea water by temperature and salinity, and the density of lake water by temperature. Where stratification occurs, there may be thin layers in which temperature or some other property changes more rapidly with height or depth than the surrounding fluid. Depending on the main sources of buoyancy, this layer may be called a pycnocline (density), thermocline (temperature), halocline (salinity), or chemocline (chemistry, including oxygenation).
The same buoyancy that gives rise to stratification also drives gravity waves. If the gravity waves occur within the fluid, they are called internal waves.
In modeling buoyancy-driven flows, the Navier-Stokes equations are modified using the Boussinesq approximation. This ignores variations in density except where they are multiplied by the gravitational acceleration .
If the pressure depends only on density and vice versa, the fluid dynamics are called barotropic. In the atmosphere, this corresponds to a lack of fronts, as in the tropics. If there are fronts, the flow is baroclinic, and instabilities such as cyclones can occur.
Rotation
Coriolis effect
Circulation
Kelvin's circulation theorem
Vorticity equation
Thermal wind
Geostrophic current
Geostrophic wind
Taylor–Proudman theorem
Hydrostatic equilibrium
Ekman spiral
Ekman layer
General circulation
Atmospheric circulation
Ocean current
Ocean dynamics
Thermohaline circulation
Boundary current
Sverdrup balance
Subsurface currents
Waves
Barotropic
Kelvin wave
Rossby wave
Sverdrup wave (Poincaré wave)
Baroclinic
Gravity wave
| Physical sciences | Geophysics | Earth science |
29007035 | https://en.wikipedia.org/wiki/Inkayacu | Inkayacu | Inkayacu is a genus of extinct penguins. It lived in what is now Peru during the Late Eocene, around 36 million years ago. A nearly complete skeleton was discovered in 2008 and includes fossilized feathers, the first known in penguins. A study of the melanosomes, pigment-containing organelles within the feathers, indicated that they were gray or reddish brown. This differs from modern penguins, which get their dark black-brown feathers from unique melanosomes that are large and ellipsoidal.
Etymology
The genus name derives from the Quechua words Inka for emperor or king and yacu for water, thus the name translates to "Water Emperor" or, more idiomatic, "king of the water". The specific epithet refers to Paracas, where the fossils were found.
Description
Although it was an early penguin, Inkayacu closely resembled its modern relatives. It had paddle-like wings with short feathers, and a very long bill. Inkayacu, along with other extinct penguins from Peru, are often referred to as giant penguins because of their large size. Inkayacu was among the largest described fossil penguins, measuring long and weighing about , twice as heavy as the average emperor penguin, the largest extant penguin.
The melanosomes within the feathers of Inkayacu are long and narrow, similar to most other birds. Their shape suggests that Inkayacu had grey and reddish-brown feathering across its body. Most modern penguins have melanosomes that are about the same length as those of Inkayacu, but are much wider. There is also a greater number of them within living penguins' cells. The shape of these melanosomes gives them a dark brown or black color, and is the reason why modern penguins are mostly black and white. Despite not having the distinctive melanosomes of modern penguins, the feathers of Inkayacu were similar in many other ways. The melanosomes within the feathers provides both color and wear resistance. The feathers that made up the body contour of the bird have large shafts, and the primaries along the edge of the wings are short and undifferentiated. The nanostructure of penguin feathers was modified after earlier macrostructural modifications of the feather shape that has been linked to aquatic flight.
Discovery
Fossils of Inkayacu were first found in 2008 on the Pacific coast of Ica, Peru. A nearly complete skeleton was uncovered from the Otuma Formation, in the Paracas National Reserve by an expedition team led by Rodolfo Salas and studied by a team led by Julia Clarke of the University of Texas. It was nicknamed, "Pedro," by the researchers that had discovered it. This had been the first recovered fossil with feathers attached to it. The feathers were well enough preserved that researchers Julia Clarke, Liliana D'Alba and Ali J. Altamirano were able to perform analysis of melanosomes. Until now, without the addition of feathers, there has not been any research conducted on the nanostructure of ancient feathers. Large penguins, including the species Perudyptes devriesi and Icadyptes salasi, had been described from the area the previous year.
The first evidence of melanosomes in fossilized feathers was published in late 2008, being reported from an Early Cretaceous bird. Paleontologist Jakob Vinther, an author of the 2007 paper on the first fossilized melanosomes known, found melanosomes in the feathers of Inkayacu soon after the fossil was discovered. This discovery gives insight into how the evolutionary history of Inkayacu has affected the morphology of its extant descendants.
Paleobiology
Inkayacu inhabited a sea that existed in Peru during the Late Eocene. Paddle-like limbs enabled an aquatic lifestyle. The large tightly packed melanosomes within the cells of living penguins gives the feathers added rigidity, which may be an adaptation for coping with the stresses of underwater flight. Because Inkayacu has smaller and fewer melanosomes, it may not have been able to swim very deep, possibly remaining near the surface. However, it is also possible that the melanosomes of modern penguins do not give them an advantage underwater, since the feathers on their undersides are primarily white, lacking the rigidity of melanin. If melanin is present in the feathers to add rigidity, it would be expected that all feathers on living penguins would be black.
| Biology and health sciences | Prehistoric birds | Animals |
2068726 | https://en.wikipedia.org/wiki/History%20of%20Earth | History of Earth | The natural history of Earth concerns the development of planet Earth from its formation to the present day. Nearly all branches of natural science have contributed to understanding of the main events of Earth's past, characterized by constant geological change and biological evolution.
The geological time scale (GTS), as defined by international convention, depicts the large spans of time from the beginning of the Earth to the present, and its divisions chronicle some definitive events of Earth history. Earth formed around 4.54 billion years ago, approximately one-third the age of the universe, by accretion from the solar nebula. Volcanic outgassing probably created the primordial atmosphere and then the ocean, but the early atmosphere contained almost no oxygen. Much of the Earth was molten because of frequent collisions with other bodies which led to extreme volcanism. While the Earth was in its earliest stage (Early Earth), a giant impact collision with a planet-sized body named Theia is thought to have formed the Moon. Over time, the Earth cooled, causing the formation of a solid crust, and allowing liquid water on the surface.
The Hadean eon represents the time before a reliable (fossil) record of life; it began with the formation of the planet and ended 4.0 billion years ago. The following Archean and Proterozoic eons produced the beginnings of life on Earth and its earliest evolution. The succeeding eon is the Phanerozoic, divided into three eras: the Palaeozoic, an era of arthropods, fishes, and the first life on land; the Mesozoic, which spanned the rise, reign, and climactic extinction of the non-avian dinosaurs; and the Cenozoic, which saw the rise of mammals. Recognizable humans emerged at most 2 million years ago, a vanishingly small period on the geological scale.
The earliest undisputed evidence of life on Earth dates at least from 3.5 billion years ago, during the Eoarchean Era, after a geological crust started to solidify following the earlier molten Hadean eon. There are microbial mat fossils such as stromatolites found in 3.48 billion-year-old sandstone discovered in Western Australia. Other early physical evidence of a biogenic substance is graphite in 3.7 billion-year-old metasedimentary rocks discovered in southwestern Greenland as well as "remains of biotic life" found in 4.1 billion-year-old rocks in Western Australia. According to one of the researchers, "If life arose relatively quickly on Earth … then it could be common in the universe."
Photosynthetic organisms appeared between 3.2 and 2.4 billion years ago and began enriching the atmosphere with oxygen. Life remained mostly small and microscopic until about 580 million years ago, when complex multicellular life arose, developed over time, and culminated in the Cambrian Explosion about 538.8 million years ago. This sudden diversification of life forms produced most of the major phyla known today, and divided the Proterozoic Eon from the Cambrian Period of the Paleozoic Era. It is estimated that 99 percent of all species that ever lived on Earth, over five billion, have gone extinct. Estimates on the number of Earth's current species range from 10 million to 14 million, of which about 1.2 million are documented, but over 86 percent have not been described.
The Earth's crust has constantly changed since its formation, as has life since its first appearance. Species continue to evolve, taking on new forms, splitting into daughter species, or going extinct in the face of ever-changing physical environments. The process of plate tectonics continues to shape the Earth's continents and oceans and the life they harbor.
Eons
In geochronology, time is generally measured in mya (million years ago), each unit representing the period of approximately 1,000,000 years in the past. The history of Earth is divided into four great eons, starting 4,540 mya with the formation of the planet. Each eon saw the most significant changes in Earth's composition, climate and life. Each eon is subsequently divided into eras, which in turn are divided into periods, which are further divided into epochs.
Geologic time scale
The history of the Earth can be organized chronologically according to the geologic time scale, which is split into intervals based on stratigraphic analysis.
Solar System formation
The standard model for the formation of the Solar System (including the Earth) is the solar nebula hypothesis. In this model, the Solar System formed from a large, rotating cloud of interstellar dust and gas called the solar nebula. It was composed of hydrogen and helium created shortly after the Big Bang 13.8 Ga (billion years ago) and heavier elements ejected by supernovae. About 4.5 Ga, the nebula began a contraction that may have been triggered by the shock wave from a nearby supernova. A shock wave would have also made the nebula rotate. As the cloud began to accelerate, its angular momentum, gravity, and inertia flattened it into a protoplanetary disk perpendicular to its axis of rotation. Small perturbations due to collisions and the angular momentum of other large debris created the means by which kilometer-sized protoplanets began to form, orbiting the nebular center.
The center of the nebula, not having much angular momentum, collapsed rapidly, the compression heating it until nuclear fusion of hydrogen into helium began. After more contraction, a T Tauri star ignited and evolved into the Sun. Meanwhile, in the outer part of the nebula gravity caused matter to condense around density perturbations and dust particles, and the rest of the protoplanetary disk began separating into rings. In a process known as runaway accretion, successively larger fragments of dust and debris clumped together to form planets. Earth formed in this manner about 4.54 billion years ago (with an uncertainty of 1%) and was largely completed within 10–20 million years. In June 2023, scientists reported evidence that the planet Earth may have formed in just three million years, much faster than the 10−100 million years thought earlier. Nonetheless, the solar wind of the newly formed T Tauri star cleared out most of the material in the disk that had not already condensed into larger bodies. The same process is expected to produce accretion disks around virtually all newly forming stars in the universe, some of which yield planets.
The proto-Earth grew by accretion until its interior was hot enough to melt the heavy, siderophile metals. Having higher densities than the silicates, these metals sank. This so-called iron catastrophe resulted in the separation of a primitive mantle and a (metallic) core only 10 million years after the Earth began to form, producing the layered structure of Earth and setting up the formation of Earth's magnetic field. J.A. Jacobs was the first to suggest that Earth's inner core—a solid center distinct from the liquid outer core—is freezing and growing out of the liquid outer core due to the gradual cooling of Earth's interior (about 100 degrees Celsius per billion years).
Hadean and Archean Eons
The first eon in Earth's history, the Hadean, begins with the Earth's formation and is followed by the Archean eon at 3.8 Ga. The oldest rocks found on Earth date to about 4.0 Ga, and the oldest detrital zircon crystals in rocks to about 4.4 Ga, soon after the formation of the Earth's crust and the Earth itself. The giant impact hypothesis for the Moon's formation states that shortly after formation of an initial crust, the proto-Earth was impacted by a smaller protoplanet, which ejected part of the mantle and crust into space and created the Moon.
From crater counts on other celestial bodies, it is inferred that a period of intense meteorite impacts, called the Late Heavy Bombardment, began about 4.1 Ga, and concluded around 3.8 Ga, at the end of the Hadean. In addition, volcanism was severe due to the large heat flow and geothermal gradient. Nevertheless, detrital zircon crystals dated to 4.4 Ga show evidence of having undergone contact with liquid water, suggesting that the Earth already had oceans or seas at that time.
By the beginning of the Archean, the Earth had cooled significantly. Present life forms could not have survived at Earth's surface, because the Archean atmosphere lacked oxygen hence had no ozone layer to block ultraviolet light. Nevertheless, it is believed that primordial life began to evolve by the early Archean, with candidate fossils dated to around 3.5 Ga. Some scientists even speculate that life could have begun during the early Hadean, as far back as 4.4 Ga, surviving the possible Late Heavy Bombardment period in hydrothermal vents below the Earth's surface.
Formation of the Moon
Earth's only natural satellite, the Moon, is larger relative to its planet than any other satellite in the Solar System. During the Apollo program, rocks from the Moon's surface were brought to Earth. Radiometric dating of these rocks shows that the Moon is 4.53 ± 0.01 billion years old, formed at least 30 million years after the Solar System. New evidence suggests the Moon formed even later, 4.48 ± 0.02 Ga, or 70–110 million years after the start of the Solar System.
Theories for the formation of the Moon must explain its late formation as well as the following facts. First, the Moon has a low density (3.3 times that of water, compared to 5.5 for the Earth) and a small metallic core. Second, the Earth and Moon have the same oxygen isotopic signature (relative abundance of the oxygen isotopes). Of the theories proposed to account for these phenomena, one is widely accepted: The giant impact hypothesis proposes that the Moon originated after a body the size of Mars (sometimes named Theia) struck the proto-Earth a glancing blow.
The collision released about 100 million times more energy than the more recent Chicxulub impact that is believed to have caused the extinction of the non-avian dinosaurs. It was enough to vaporize some of the Earth's outer layers and melt both bodies. A portion of the mantle material was ejected into orbit around the Earth. The giant impact hypothesis predicts that the Moon was depleted of metallic material, explaining its abnormal composition. The ejecta in orbit around the Earth could have condensed into a single body within a couple of weeks. Under the influence of its own gravity, the ejected material became a more spherical body: the Moon.
First continents
Mantle convection, the process that drives plate tectonics, is a result of heat flow from the Earth's interior to the Earth's surface. It involves the creation of rigid tectonic plates at mid-oceanic ridges. These plates are destroyed by subduction into the mantle at subduction zones. During the early Archean (about 3.0 Ga) the mantle was much hotter than today, probably around , so convection in the mantle was faster. Although a process similar to present-day plate tectonics did occur, this would have gone faster too. It is likely that during the Hadean and Archean, subduction zones were more common, and therefore tectonic plates were smaller.
The initial crust, which formed when the Earth's surface first solidified, totally disappeared from a combination of this fast Hadean plate tectonics and the intense impacts of the Late Heavy Bombardment. However, it is thought that it was basaltic in composition, like today's oceanic crust, because little crustal differentiation had yet taken place. The first larger pieces of continental crust, which is a product of differentiation of lighter elements during partial melting in the lower crust, appeared at the end of the Hadean, about 4.0 Ga. What is left of these first small continents are called cratons. These pieces of late Hadean and early Archean crust form the cores around which today's continents grew.
The oldest rocks on Earth are found in the North American craton of Canada. They are tonalites from about 4.0 Ga. They show traces of metamorphism by high temperature, but also sedimentary grains that have been rounded by erosion during transport by water, showing that rivers and seas existed then. Cratons consist primarily of two alternating types of terranes. The first are so-called greenstone belts, consisting of low-grade metamorphosed sedimentary rocks. These "greenstones" are similar to the sediments today found in oceanic trenches, above subduction zones. For this reason, greenstones are sometimes seen as evidence for subduction during the Archean. The second type is a complex of felsic magmatic rocks. These rocks are mostly tonalite, trondhjemite or granodiorite, types of rock similar in composition to granite (hence such terranes are called TTG-terranes). TTG-complexes are seen as the relicts of the first continental crust, formed by partial melting in basalt.
Oceans and atmosphere
Earth is often described as having had three atmospheres. The first atmosphere, captured from the solar nebula, was composed of light (atmophile) elements from the solar nebula, mostly hydrogen and helium. A combination of the solar wind and Earth's heat would have driven off this atmosphere, as a result of which the atmosphere is now depleted of these elements compared to cosmic abundances. After the impact which created the Moon, the molten Earth released volatile gases; and later more gases were released by volcanoes, completing a second atmosphere rich in greenhouse gases but poor in oxygen. Finally, the third atmosphere, rich in oxygen, emerged when bacteria began to produce oxygen about 2.8 Ga.
In early models for the formation of the atmosphere and ocean, the second atmosphere was formed by outgassing of volatiles from the Earth's interior. Now it is considered likely that many of the volatiles were delivered during accretion by a process known as impact degassing in which incoming bodies vaporize on impact. The ocean and atmosphere would, therefore, have started to form even as the Earth formed. The new atmosphere probably contained water vapor, carbon dioxide, nitrogen, and smaller amounts of other gases.
Planetesimals at a distance of 1 astronomical unit (AU), the distance of the Earth from the Sun, probably did not contribute any water to the Earth because the solar nebula was too hot for ice to form and the hydration of rocks by water vapor would have taken too long. The water must have been supplied by meteorites from the outer asteroid belt and some large planetary embryos from beyond 2.5 AU. Comets may also have contributed. Though most comets are today in orbits farther away from the Sun than Neptune, computer simulations show that they were originally far more common in the inner parts of the Solar System.
As the Earth cooled, clouds formed. Rain created the oceans. Recent evidence suggests the oceans may have begun forming as early as 4.4 Ga. By the start of the Archean eon, they already covered much of the Earth. This early formation has been difficult to explain because of a problem known as the faint young Sun paradox. Stars are known to get brighter as they age, and the Sun has become 30% brighter since its formation 4.5 billion years ago. Many models indicate that the early Earth should have been covered in ice. A likely solution is that there was enough carbon dioxide and methane to produce a greenhouse effect. The carbon dioxide would have been produced by volcanoes and the methane by early microbes. It is hypothesized that there also existed an organic haze created from the products of methane photolysis that caused an anti-greenhouse effect as well. Another greenhouse gas, ammonia, would have been ejected by volcanos but quickly destroyed by ultraviolet radiation.
Origin of life
One of the reasons for interest in the early atmosphere and ocean is that they form the conditions under which life first arose. There are many models, but little consensus, on how life emerged from non-living chemicals; chemical systems created in the laboratory fall well short of the minimum complexity for a living organism.
The first step in the emergence of life may have been chemical reactions that produced many of the simpler organic compounds, including nucleobases and amino acids, that are the building blocks of life. An experiment in 1952 by Stanley Miller and Harold Urey showed that such molecules could form in an atmosphere of water, methane, ammonia and hydrogen with the aid of sparks to mimic the effect of lightning. Although atmospheric composition was probably different from that used by Miller and Urey, later experiments with more realistic compositions also managed to synthesize organic molecules. Computer simulations show that extraterrestrial organic molecules could have formed in the protoplanetary disk before the formation of the Earth.
Additional complexity could have been reached from at least three possible starting points: self-replication, an organism's ability to produce offspring that are similar to itself; metabolism, its ability to feed and repair itself; and external cell membranes, which allow food to enter and waste products to leave, but exclude unwanted substances.
Replication first: RNA world
Even the simplest members of the three modern domains of life use DNA to record their "recipes" and a complex array of RNA and protein molecules to "read" these instructions and use them for growth, maintenance, and self-replication.
The discovery that a kind of RNA molecule called a ribozyme can catalyze both its own replication and the construction of proteins led to the hypothesis that earlier life-forms were based entirely on RNA. They could have formed an RNA world in which there were individuals but no species, as mutations and horizontal gene transfers would have meant that the offspring in each generation were quite likely to have different genomes from those that their parents started with. RNA would later have been replaced by DNA, which is more stable and therefore can build longer genomes, expanding the range of capabilities a single organism can have. Ribozymes remain as the main components of ribosomes, the "protein factories" of modern cells.
Although short, self-replicating RNA molecules have been artificially produced in laboratories, doubts have been raised about whether natural non-biological synthesis of RNA is possible. The earliest ribozymes may have been formed of simpler nucleic acids such as PNA, TNA or GNA, which would have been replaced later by RNA. Other pre-RNA replicators have been posited, including crystals and even quantum systems.
In 2003 it was proposed that porous metal sulfide precipitates would assist RNA synthesis at about and at ocean-bottom pressures near hydrothermal vents. In this hypothesis, the proto-cells would be confined in the pores of the metal substrate until the later development of lipid membranes.
Metabolism first: iron–sulfur world
Another long-standing hypothesis is that the first life was composed of protein molecules. Amino acids, the building blocks of proteins, are easily synthesized in plausible prebiotic conditions, as are small peptides (polymers of amino acids) that make good catalysts. A series of experiments starting in 1997 showed that amino acids and peptides could form in the presence of carbon monoxide and hydrogen sulfide with iron sulfide and nickel sulfide as catalysts. Most of the steps in their assembly required temperatures of about and moderate pressures, although one stage required and a pressure equivalent to that found under of rock. Hence, self-sustaining synthesis of proteins could have occurred near hydrothermal vents.
A difficulty with the metabolism-first scenario is finding a way for organisms to evolve. Without the ability to replicate as individuals, aggregates of molecules would have "compositional genomes" (counts of molecular species in the aggregate) as the target of natural selection. However, a recent model shows that such a system is unable to evolve in response to natural selection.
Membranes first: Lipid world
It has been suggested that double-walled "bubbles" of lipids like those that form the external membranes of cells may have been an essential first step. Experiments that simulated the conditions of the early Earth have reported the formation of lipids, and these can spontaneously form liposomes, double-walled "bubbles", and then reproduce themselves. Although they are not intrinsically information-carriers as nucleic acids are, they would be subject to natural selection for longevity and reproduction. Nucleic acids such as RNA might then have formed more easily within the liposomes than they would have outside.
The clay theory
Some clays, notably montmorillonite, have properties that make them plausible accelerators for the emergence of an RNA world: they grow by self-replication of their crystalline pattern, are subject to an analog of natural selection (as the clay "species" that grows fastest in a particular environment rapidly becomes dominant), and can catalyze the formation of RNA molecules. Although this idea has not become the scientific consensus, it still has active supporters.
Research in 2003 reported that montmorillonite could also accelerate the conversion of fatty acids into "bubbles", and that the bubbles could encapsulate RNA attached to the clay. Bubbles can then grow by absorbing additional lipids and dividing. The formation of the earliest cells may have been aided by similar processes.
A similar hypothesis presents self-replicating iron-rich clays as the progenitors of nucleotides, lipids and amino acids.
Last universal common ancestor
It is believed that of this multiplicity of protocells, only one line survived. Current phylogenetic evidence suggests that the last universal ancestor (LUA) lived during the early Archean eon, perhaps 3.5 Ga or earlier. This LUA cell is the ancestor of all life on Earth today. It was probably a prokaryote, possessing a cell membrane and probably ribosomes, but lacking a nucleus or membrane-bound organelles such as mitochondria or chloroplasts. Like modern cells, it used DNA as its genetic code, RNA for information transfer and protein synthesis, and enzymes to catalyze reactions. Some scientists believe that instead of a single organism being the last universal common ancestor, there were populations of organisms exchanging genes by lateral gene transfer.
Proterozoic Eon
The Proterozoic eon lasted from 2.5 Ga to 538.8 Ma (million years) ago. In this time span, cratons grew into continents with modern sizes. The change to an oxygen-rich atmosphere was a crucial development. Life developed from prokaryotes into eukaryotes and multicellular forms. The Proterozoic saw a couple of severe ice ages called Snowball Earths. After the last Snowball Earth about 600 Ma, the evolution of life on Earth accelerated. About 580 Ma, the Ediacaran biota formed the prelude for the Cambrian Explosion.
Oxygen revolution
The earliest cells absorbed energy and food from the surrounding environment. They used fermentation, the breakdown of more complex compounds into less complex compounds with less energy, and used the energy so liberated to grow and reproduce. Fermentation can only occur in an anaerobic (oxygen-free) environment. The evolution of photosynthesis made it possible for cells to derive energy from the Sun.
Most of the life that covers the surface of the Earth depends directly or indirectly on photosynthesis. The most common form, oxygenic photosynthesis, turns carbon dioxide, water, and sunlight into food. It captures the energy of sunlight in energy-rich molecules such as ATP, which then provide the energy to make sugars. To supply the electrons in the circuit, hydrogen is stripped from water, leaving oxygen as a waste product. Some organisms, including purple bacteria and green sulfur bacteria, use an anoxygenic form of photosynthesis that uses alternatives to hydrogen stripped from water as electron donors; examples are hydrogen sulfide, sulfur and iron. Such extremophile organisms are restricted to otherwise inhospitable environments such as hot springs and hydrothermal vents.
The simpler anoxygenic form arose about 3.8 Ga, not long after the appearance of life. The timing of oxygenic photosynthesis is more controversial; it had certainly appeared by about 2.4 Ga, but some researchers put it back as far as 3.2 Ga. The latter "probably increased global productivity by at least two or three orders of magnitude". Among the oldest remnants of oxygen-producing lifeforms are fossil stromatolites.
At first, the released oxygen was bound up with limestone, iron, and other minerals. The oxidized iron appears as red layers in geological strata called banded iron formations that formed in abundance during the Siderian period (between 2500 Ma and 2300 Ma). When most of the exposed readily reacting minerals were oxidized, oxygen finally began to accumulate in the atmosphere. Though each cell only produced a minute amount of oxygen, the combined metabolism of many cells over a vast time transformed Earth's atmosphere to its current state. This was Earth's third atmosphere.
Some oxygen was stimulated by solar ultraviolet radiation to form ozone, which collected in a layer near the upper part of the atmosphere. The ozone layer absorbed, and still absorbs, a significant amount of the ultraviolet radiation that once had passed through the atmosphere. It allowed cells to colonize the surface of the ocean and eventually the land: without the ozone layer, ultraviolet radiation bombarding land and sea would have caused unsustainable levels of mutation in exposed cells.
Photosynthesis had another major impact. Oxygen was toxic; much life on Earth probably died out as its levels rose in what is known as the oxygen catastrophe. Resistant forms survived and thrived, and some developed the ability to use oxygen to increase their metabolism and obtain more energy from the same food.
Snowball Earth
The natural evolution of the Sun made it progressively more luminous during the Archean and Proterozoic eons; the Sun's luminosity increases 6% every billion years. As a result, the Earth began to receive more heat from the Sun in the Proterozoic eon. However, the Earth did not get warmer. Instead, the geological record suggests it cooled dramatically during the early Proterozoic. Glacial deposits found in South Africa date back to 2.2 Ga, at which time, based on paleomagnetic evidence, they must have been located near the equator. Thus, this glaciation, known as the Huronian glaciation, may have been global. Some scientists suggest this was so severe that the Earth was frozen over from the poles to the equator, a hypothesis called Snowball Earth.
The Huronian ice age might have been caused by the increased oxygen concentration in the atmosphere, which caused the decrease of methane (CH4) in the atmosphere. Methane is a strong greenhouse gas, but with oxygen it reacts to form CO2, a less effective greenhouse gas. When free oxygen became available in the atmosphere, the concentration of methane could have decreased dramatically, enough to counter the effect of the increasing heat flow from the Sun.
However, the term Snowball Earth is more commonly used to describe later extreme ice ages during the Cryogenian period. There were four periods, each lasting about 10 million years, between 750 and 580 million years ago, when the Earth is thought to have been covered with ice apart from the highest mountains, and average temperatures were about . The snowball may have been partly due to the location of the supercontinent Rodinia straddling the Equator. Carbon dioxide combines with rain to weather rocks to form carbonic acid, which is then washed out to sea, thus extracting the greenhouse gas from the atmosphere. When the continents are near the poles, the advance of ice covers the rocks, slowing the reduction in carbon dioxide, but in the Cryogenian the weathering of Rodinia was able to continue unchecked until the ice advanced to the tropics. The process may have finally been reversed by the emission of carbon dioxide from volcanoes or the destabilization of methane gas hydrates. According to the alternative Slushball Earth theory, even at the height of the ice ages there was still open water at the Equator.
Emergence of eukaryotes
Modern taxonomy classifies life into three domains. The time of their origin is uncertain. The Bacteria domain probably first split off from the other forms of life (sometimes called Neomura), but this supposition is controversial. Soon after this, by 2 Ga, the Neomura split into the Archaea and the Eukaryota. Eukaryotic cells (Eukaryota) are larger and more complex than prokaryotic cells (Bacteria and Archaea), and the origin of that complexity is only now becoming known. The earliest fossils possessing features typical of fungi date to the Paleoproterozoic era, some 2.4 Ga ago; these multicellular benthic organisms had filamentous structures capable of anastomosis.
Around this time, the first proto-mitochondrion was formed. A bacterial cell related to today's Rickettsia, which had evolved to metabolize oxygen, entered a larger prokaryotic cell, which lacked that capability. Perhaps the large cell attempted to digest the smaller one but failed (possibly due to the evolution of prey defenses). The smaller cell may have tried to parasitize the larger one. In any case, the smaller cell survived inside the larger cell. Using oxygen, it metabolized the larger cell's waste products and derived more energy. Part of this excess energy was returned to the host. The smaller cell replicated inside the larger one. Soon, a stable symbiosis developed between the large cell and the smaller cells inside it. Over time, the host cell acquired some genes from the smaller cells, and the two kinds became dependent on each other: the larger cell could not survive without the energy produced by the smaller ones, and these, in turn, could not survive without the raw materials provided by the larger cell. The whole cell is now considered a single organism, and the smaller cells are classified as organelles called mitochondria.
A similar event occurred with photosynthetic cyanobacteria entering large heterotrophic cells and becoming chloroplasts. Probably as a result of these changes, a line of cells capable of photosynthesis split off from the other eukaryotes more than 1 billion years ago. There were probably several such inclusion events. Besides the well-established endosymbiotic theory of the cellular origin of mitochondria and chloroplasts, there are theories that cells led to peroxisomes, spirochetes led to cilia and flagella, and that perhaps a DNA virus led to the cell nucleus, though none of them are widely accepted.
Archaeans, bacteria, and eukaryotes continued to diversify and to become more complex and better adapted to their environments. Each domain repeatedly split into multiple lineages. Around 1.1 Ga, the plant, animal, and fungi lines had split, though they still existed as solitary cells. Some of these lived in colonies, and gradually a division of labor began to take place; for instance, cells on the periphery might have started to assume different roles from those in the interior. Although the division between a colony with specialized cells and a multicellular organism is not always clear, around 1 billion years ago, the first multicellular plants emerged, probably green algae. Possibly by around 900 Ma true multicellularity had also evolved in animals.
At first, it probably resembled today's sponges, which have totipotent cells that allow a disrupted organism to reassemble itself. As the division of labor was completed in the different lineages of multicellular organisms, cells became more specialized and more dependent on each other.
Supercontinents in the Proterozoic
Reconstructions of tectonic plate movement in the past 250 million years (the Cenozoic and Mesozoic eras) can be made reliably using fitting of continental margins, ocean floor magnetic anomalies and paleomagnetic poles. No ocean crust dates back further than that, so earlier reconstructions are more difficult. Paleomagnetic poles are supplemented by geologic evidence such as orogenic belts, which mark the edges of ancient plates, and past distributions of flora and fauna. The further back in time, the scarcer and harder to interpret the data get and the more uncertain the reconstructions.
Throughout the history of the Earth, there have been times when continents collided and formed a supercontinent, which later broke up into new continents. About 1000 to 830 Ma, most continental mass was united in the supercontinent Rodinia. Rodinia may have been preceded by Early-Middle Proterozoic continents called Nuna and Columbia.
After the break-up of Rodinia about 800 Ma, the continents may have formed another short-lived supercontinent around 550 Ma. The hypothetical supercontinent is sometimes referred to as Pannotia or Vendia. The evidence for it is a phase of continental collision known as the Pan-African orogeny, which joined the continental masses of current-day Africa, South America, Antarctica and Australia. The existence of Pannotia depends on the timing of the rifting between Gondwana (which included most of the landmass now in the Southern Hemisphere, as well as the Arabian Peninsula and the Indian subcontinent) and Laurentia (roughly equivalent to current-day North America). It is at least certain that by the end of the Proterozoic eon, most of the continental mass lay united in a position around the south pole.
Late Proterozoic climate and life
The end of the Proterozoic saw at least two Snowball Earths, so severe that the surface of the oceans may have been completely frozen. This happened about 716.5 and 635 Ma, in the Cryogenian period. The intensity and mechanism of both glaciations are still under investigation and harder to explain than the early Proterozoic Snowball Earth.
Most paleoclimatologists think the cold episodes were linked to the formation of the supercontinent Rodinia. Because Rodinia was centered on the equator, rates of chemical weathering increased and carbon dioxide (CO2) was taken from the atmosphere. Because CO2 is an important greenhouse gas, climates cooled globally.
In the same way, during the Snowball Earths most of the continental surface was covered with permafrost, which decreased chemical weathering again, leading to the end of the glaciations. An alternative hypothesis is that enough carbon dioxide escaped through volcanic outgassing that the resulting greenhouse effect raised global temperatures. Increased volcanic activity resulted from the break-up of Rodinia at about the same time.
The Cryogenian period was followed by the Ediacaran period, which was characterized by a rapid development of new multicellular lifeforms. Whether there is a connection between the end of the severe ice ages and the increase in diversity of life is not clear, but it does not seem coincidental. The new forms of life, called Ediacara biota, were larger and more diverse than ever. Though the taxonomy of most Ediacaran life forms is unclear, some were ancestors of groups of modern life. Important developments were the origin of muscular and neural cells. None of the Ediacaran fossils had hard body parts like skeletons. These first appear after the boundary between the Proterozoic and Phanerozoic eons or Ediacaran and Cambrian periods.
Phanerozoic Eon
The Phanerozoic is the current eon on Earth, which started approximately 538.8 million years ago. It consists of three eras: The Paleozoic, Mesozoic, and Cenozoic, and is the time when multi-cellular life greatly diversified into almost all the organisms known today.
The Paleozoic ("old life") era was the first and longest era of the Phanerozoic eon, lasting from 538.8 to 251.9 Ma. During the Paleozoic, many modern groups of life came into existence. Life colonized the land, first plants, then animals. Two significant extinctions occurred. The continents formed at the break-up of Pannotia and Rodinia at the end of the Proterozoic slowly moved together again, forming the supercontinent Pangaea in the late Paleozoic.
The Mesozoic ("middle life") era lasted from 251.9 Ma to 66 Ma. It is subdivided into the Triassic, Jurassic, and Cretaceous periods. The era began with the Permian–Triassic extinction event, the most severe extinction event in the fossil record; 95% of the species on Earth died out. It ended with the Cretaceous–Paleogene extinction event that wiped out the dinosaurs.
The Cenozoic ("new life") era began at Ma, and is subdivided into the Paleogene, Neogene, and Quaternary periods. These three periods are further split into seven subdivisions, with the Paleogene composed of The Paleocene, Eocene, and Oligocene, the Neogene divided into the Miocene, Pliocene, and the Quaternary composed of the Pleistocene, and Holocene. Mammals, birds, amphibians, crocodilians, turtles, and lepidosaurs survived the Cretaceous–Paleogene extinction event that killed off the non-avian dinosaurs and many other forms of life, and this is the era during which they diversified into their modern forms.
Tectonics, paleogeography and climate
At the end of the Proterozoic, the supercontinent Pannotia had broken apart into the smaller continents Laurentia, Baltica, Siberia and Gondwana. During periods when continents move apart, more oceanic crust is formed by volcanic activity. Because the young volcanic crust is relatively hotter and less dense than the old oceanic crust, the ocean floors rise during such periods. This causes the sea level to rise. Therefore, in the first half of the Paleozoic, large areas of the continents were below sea level.
Early Paleozoic climates were warmer than today, but the end of the Ordovician saw a short ice age during which glaciers covered the south pole, where the huge continent Gondwana was situated. Traces of glaciation from this period are only found on former Gondwana. During the Late Ordovician ice age, a few mass extinctions took place, in which many brachiopods, trilobites, Bryozoa and corals disappeared. These marine species could probably not contend with the decreasing temperature of the sea water.
The continents Laurentia and Baltica collided between 450 and 400 Ma, during the Caledonian Orogeny, to form Laurussia (also known as Euramerica). Traces of the mountain belt this collision caused can be found in Scandinavia, Scotland, and the northern Appalachians. In the Devonian period (416–359 Ma) Gondwana and Siberia began to move towards Laurussia. The collision of Siberia with Laurussia caused the Uralian Orogeny, the collision of Gondwana with Laurussia is called the Variscan or Hercynian Orogeny in Europe or the Alleghenian Orogeny in North America. The latter phase took place during the Carboniferous period (359–299 Ma) and resulted in the formation of the last supercontinent, Pangaea.
By 180 Ma, Pangaea broke up into Laurasia and Gondwana.
Cambrian explosion
The rate of the evolution of life as recorded by fossils accelerated in the Cambrian period (542–488 Ma). The sudden emergence of many new species, phyla, and forms in this period is called the Cambrian Explosion. It was a form of adaptive radiation, where vacant niches left by the extinct Ediacaran biota were filled up by the emergence of new phyla. The biological fomenting in the Cambrian Explosion was unprecedented before and since that time. Whereas the Ediacaran life forms appear yet primitive and not easy to put in any modern group, at the end of the Cambrian, most modern phyla were already present. The development of hard body parts such as shells, skeletons or exoskeletons in animals like molluscs, echinoderms, crinoids and arthropods (a well-known group of arthropods from the lower Paleozoic are the trilobites) made the preservation and fossilization of such life forms easier than those of their Proterozoic ancestors. For this reason, much more is known about life in and after the Cambrian period than about life in older periods. Some of these Cambrian groups appear complex but are seemingly quite different from modern life; examples are Anomalocaris and Haikouichthys. More recently, however, these seem to have found a place in modern classification.
During the Cambrian, the first vertebrate animals, among them the first fishes, had appeared. A creature that could have been the ancestor of the fishes, or was probably closely related to it, was Pikaia. It had a primitive notochord, a structure that could have developed into a vertebral column later. The first fishes with jaws (Gnathostomata) appeared during the next geological period, the Ordovician. The colonisation of new niches resulted in massive body sizes. In this way, fishes with increasing sizes evolved during the early Paleozoic, such as the titanic placoderm Dunkleosteus, which could grow long.
The diversity of life forms did not increase significantly because of a series of mass extinctions that define widespread biostratigraphic units called biomeres. After each extinction pulse, the continental shelf regions were repopulated by similar life forms that may have been evolving slowly elsewhere. By the late Cambrian, the trilobites had reached their greatest diversity and dominated nearly all fossil assemblages.
Colonization of land
Oxygen accumulation from photosynthesis resulted in the formation of an ozone layer that absorbed much of the Sun's ultraviolet radiation, meaning unicellular organisms that reached land were less likely to die, and prokaryotes began to multiply and become better adapted to survival out of the water. Prokaryote lineages had probably colonized the land as early as 3 Ga even before the origin of the eukaryotes. For a long time, the land remained barren of multicellular organisms. The supercontinent Pannotia formed around 600 Ma and then broke apart a short 50 million years later. Fish, the earliest vertebrates, evolved in the oceans around 530 Ma. A major extinction event occurred near the end of the Cambrian period, which ended 488 Ma.
Several hundred million years ago, plants (probably resembling algae) and fungi started growing at the edges of the water and then out of it. The oldest fossils of land fungi and plants date to 480–460 Ma, though molecular evidence suggests the fungi may have colonized the land as early as 1000 Ma and the plants 700 Ma. Initially remaining close to the water's edge, mutations and variations resulted in further colonization of this new environment. The timing of the first animals to leave the oceans is not precisely known: the oldest clear evidence is of arthropods on land around 450 Ma, perhaps thriving and becoming better adapted due to the vast food source provided by the terrestrial plants. There is also unconfirmed evidence that arthropods may have appeared on land as early as 530 Ma.
Evolution of tetrapods
At the end of the Ordovician period, 443 Ma, additional extinction events occurred, perhaps due to a concurrent ice age. Around 380 to 375 Ma, the first tetrapods evolved from fish. Fins evolved to become limbs that the first tetrapods used to lift their heads out of the water to breathe air. This would let them live in oxygen-poor water, or pursue small prey in shallow water. They may have later ventured on land for brief periods. Eventually, some of them became so well adapted to terrestrial life that they spent their adult lives on land, although they hatched in the water and returned to lay their eggs. This was the origin of the amphibians. About 365 Ma, another period of extinction occurred, perhaps as a result of global cooling. Plants evolved seeds, which dramatically accelerated their spread on land, around this time (by approximately 360 Ma).
About 20 million years later (340 Ma), the amniotic egg evolved, which could be laid on land, giving a survival advantage to tetrapod embryos. This resulted in the divergence of amniotes from amphibians. Another 30 million years (310 Ma) saw the divergence of the synapsids (including mammals) from the sauropsids (including birds and reptiles). Other groups of organisms continued to evolve, and lines diverged—in fish, insects, bacteria, and so on—but less is known of the details.
After yet another, the most severe extinction of the period (251~250 Ma), around 230 Ma, dinosaurs split off from their reptilian ancestors. The Triassic–Jurassic extinction event at 200 Ma spared many of the dinosaurs, and they soon became dominant among the vertebrates. Though some mammalian lines began to separate during this period, existing mammals were probably small animals resembling shrews.
The boundary between avian and non-avian dinosaurs is unclear, but Archaeopteryx, traditionally considered one of the first birds, lived around 150 Ma.
The earliest evidence for the angiosperms evolving flowers is during the Cretaceous period, some 20 million years later (132 Ma).
Extinctions
The first of five great mass extinctions was the Ordovician-Silurian extinction. Its possible cause was the intense glaciation of Gondwana, which eventually led to a Snowball Earth. 60% of marine invertebrates became extinct, and 25% of all families.
The second mass extinction was the Late Devonian extinction, probably caused by the evolution of trees, which could have led to the depletion of greenhouse gases (like ) or the eutrophication of water. 70% of all species became extinct.
The third mass extinction was the Permian-Triassic, or the Great Dying, event. The event was possibly caused by some combination of the Siberian Traps volcanic event, an asteroid impact, methane hydrate gasification, sea level fluctuations, and a major anoxic event. Either the proposed Wilkes Land crater in Antarctica or Bedout structure off the northwest coast of Australia may indicate an impact connection with the Permian-Triassic extinction. But it remains uncertain whether these or other proposed Permian-Triassic boundary craters are real impact craters or even contemporary with the Permian-Triassic extinction event. This was by far the deadliest extinction ever, with about 57% of all families and 83% of all genera killed.
The fourth mass extinction was the Triassic-Jurassic extinction event in which almost all synapsids and archosaurs became extinct, probably due to new competition from dinosaurs.
The fifth and most recent mass extinction was the Cretaceous-Paleogene extinction event. In 66 Ma, a asteroid struck Earth just off the Yucatán Peninsula—somewhere in the southwestern tip of then Laurasia—where the Chicxulub crater is today. This ejected vast quantities of particulate matter and vapor into the air that occluded sunlight, inhibiting photosynthesis. 75% of all life, including the non-avian dinosaurs, became extinct, marking the end of the Cretaceous period and Mesozoic era.
Diversification of mammals
The first true mammals evolved in the shadows of dinosaurs and other large archosaurs that filled the world by the late Triassic. The first mammals were very small, and were probably nocturnal to escape predation. Mammal diversification truly began only after the Cretaceous-Paleogene extinction event. By the early Paleocene the Earth recovered from the extinction, and mammalian diversity increased. Creatures like Ambulocetus took to the oceans to eventually evolve into whales, whereas some creatures, like primates, took to the trees. This all changed during the mid to late Eocene when the circum-Antarctic current formed between Antarctica and Australia which disrupted weather patterns on a global scale. Grassless savanna began to predominate much of the landscape, and mammals such as Andrewsarchus rose up to become the largest known terrestrial predatory mammal ever, and early whales like Basilosaurus took control of the seas.
The evolution of grasses brought a remarkable change to the Earth's landscape, and the new open spaces created pushed mammals to get bigger and bigger. Grass started to expand in the Miocene, and the Miocene is where many modern- day mammals first appeared. Giant ungulates like Paraceratherium and Deinotherium evolved to rule the grasslands. The evolution of grass also brought primates down from the trees, and started human evolution. The first big cats evolved during this time as well. The Tethys Sea was closed off by the collision of Africa and Europe.
The formation of Panama was perhaps the most important geological event to occur in the last 60 million years. Atlantic and Pacific currents were closed off from each other, which caused the formation of the Gulf Stream, which made Europe warmer. The land bridge allowed the isolated creatures of South America to migrate over to North America and vice versa. Various species migrated south, leading to the presence in South America of llamas, the spectacled bear, kinkajous and jaguars.
Three million years ago saw the start of the Pleistocene epoch, which featured dramatic climatic changes due to the ice ages. The ice ages led to the evolution and expansion of modern man in Saharan Africa. The mega-fauna that dominated fed on grasslands that, by now, had taken over much of the subtropical world. The large amounts of water held in the ice allowed various water bodies to shrink and sometimes disappear, such as the North Sea and the Bering Strait. It is believed by many that a huge migration took place along Beringia, which is why, today, there are camels (which evolved and became extinct in North America), horses (which evolved and became extinct in North America), and Native Americans. The end of the last ice age coincided with the expansion of man and a massive die out of ice age mega-fauna. This extinction is nicknamed "the Sixth Extinction".
Human evolution
A small African ape living around 6 Ma was the last animal whose descendants would include both modern humans and their closest relatives, the chimpanzees. Only two branches of its family tree have surviving descendants. Very soon after the split, for reasons that are still unclear, apes in one branch developed the ability to walk upright. Brain size increased rapidly, and by 2 Ma, the first animals classified in the genus Homo had appeared. Around the same time, the other branch split into the ancestors of the common chimpanzee and the ancestors of the bonobo as evolution continued simultaneously in all life forms.
The ability to control fire probably began in Homo erectus (or Homo ergaster), probably at least 790,000 years ago but perhaps as early as 1.5 Ma. The use and discovery of controlled fire may even predate Homo erectus. Fire was possibly used by the early Lower Paleolithic (Oldowan) hominid Homo habilis or strong australopithecines such as Paranthropus.
It is more difficult to establish the origin of language; it is unclear whether Homo erectus could speak or if that capability had not begun until Homo sapiens. As brain size increased, babies were born earlier, before their heads grew too large to pass through the pelvis. As a result, they exhibited more plasticity, thus possessing an increased capacity to learn and requiring a longer period of dependence. Social skills became more complex, language became more sophisticated, and tools became more elaborate. This contributed to further cooperation and intellectual development. Modern humans (Homo sapiens) are believed to have originated around 200,000 years ago or earlier in Africa; the oldest fossils date back to around 160,000 years ago.
The first humans to show signs of spirituality are the Neanderthals (usually classified as a separate species with no surviving descendants); they buried their dead, often with no sign of food or tools. However, evidence of more sophisticated beliefs, such as the early Cro-Magnon cave paintings (probably with magical or religious significance) did not appear until 32,000 years ago. Cro-Magnons also left behind stone figurines such as Venus of Willendorf, probably also signifying religious belief. By 11,000 years ago, Homo sapiens had reached the southern tip of South America, the last of the uninhabited continents (except for Antarctica, which remained undiscovered until 1820 AD). Tool use and communication continued to improve, and interpersonal relationships became more intricate.
Human history
Throughout more than 90% of its history, Homo sapiens lived in small bands as nomadic hunter-gatherers. As language became more complex, the ability to remember and communicate information resulted in a new replicator: the meme. Ideas could be exchanged quickly and passed down the generations. Cultural evolution quickly outpaced biological evolution, and history proper began. Between 8500 and 7000 BC, humans in the Fertile Crescent in the Middle East began the systematic husbandry of plants and animals: agriculture. This spread to neighboring regions and developed independently elsewhere until most Homo sapiens lived sedentary lives in permanent settlements as farmers. Not all societies abandoned nomadism, especially those in isolated areas of the globe poor in domesticable plant species, such as Australia. However, among those civilizations that did adopt agriculture, the relative stability and increased productivity provided by farming allowed the population to expand.
Agriculture had a major impact; humans began to affect the environment as never before. Surplus food allowed a priestly or governing class to arise, followed by increasing division of labor. This led to Earth's first civilization at Sumer in the Middle East, between 4000 and 3000 BC. Additional civilizations quickly arose in ancient Egypt, at the Indus River valley and in China. The invention of writing enabled complex societies to arise: record-keeping and libraries served as a storehouse of knowledge and increased the cultural transmission of information. Humans no longer had to spend all their time working for survival, enabling the first specialized occupations (e.g. craftsmen, merchants, priests, etc.). Curiosity and education drove the pursuit of knowledge and wisdom, and various disciplines, including science (in a primitive form), arose. This in turn led to the emergence of increasingly larger and more complex civilizations, such as the first empires, which at times traded with one another, or fought for territory and resources.
By around 500 BC, there were advanced civilizations in the Middle East, Iran, India, China, and Greece, at times expanding, at times entering into decline. In 221 BC, China became a single polity that would grow to spread its culture throughout East Asia, and it has remained the most populous nation in the world. During this period, famous Hindu texts known as vedas came in existence in Indus valley civilization. This civilization developed in warfare, arts, science, mathematics and architecture. The fundamentals of Western civilization were largely shaped in Ancient Greece, with the world's first democratic government and major advances in philosophy and science, and in Ancient Rome with advances in law, government, and engineering. The Roman Empire was Christianized by Emperor Constantine in the early 4th century and declined by the end of the 5th. Beginning with the 7th century, Christianization of Europe began, and since at least the 4th century Christianity has played a prominent role in the shaping of Western civilization. In 610, Islam was founded and quickly became the dominant religion in Western Asia. The House of Wisdom was established in Abbasid-era Baghdad, Iraq. It is considered to have been a major intellectual center during the Islamic Golden Age, where Muslim scholars in Baghdad and Cairo flourished from the ninth to the thirteenth centuries until the Mongol sack of Baghdad in 1258 AD. In 1054 AD the Great Schism between the Roman Catholic Church and the Eastern Orthodox Church led to the prominent cultural differences between Western and Eastern Europe.
In the 14th century, the Renaissance began in Italy with advances in religion, art, and science. At that time the Christian Church as a political entity lost much of its power. In 1492, Christopher Columbus reached the Americas, initiating great changes to the new world. European civilization began to change beginning in 1500, leading to the scientific and industrial revolutions. That continent began to exert political and cultural dominance over human societies around the world, a time known as the Colonial era (also see Age of Discovery). In the 18th century a cultural movement known as the Age of Enlightenment further shaped the mentality of Europe and contributed to its secularization.
| Physical sciences | Earth science | null |
2070264 | https://en.wikipedia.org/wiki/Calamine | Calamine | Calamine, also known as calamine lotion, is a medication made from powdered calamine mineral that is used to treat mild itchiness. Conditions treated include sunburn, insect bites, poison ivy, poison oak, and other mild skin conditions. It may also help dry out secretions resulting from skin irritation. It is applied on the skin as a cream or lotion.
Side effects may include skin irritation. It is considered to be safe in pregnancy. Calamine is a combination of zinc oxide and 0.5% ferric oxide (Fe2O3). The lotion is produced with additional ingredients such as phenol and calcium hydroxide.
The use of calamine lotion dates back as far as 1500 BC. It is on the World Health Organization's List of Essential Medicines. Calamine is available over-the-counter as a generic medication.
Medical uses
Calamine is used to treat itchiness. This includes sunburn, insect bite, or other mild skin conditions.
Effectiveness
The FDA recommends applying some topical over-the-counter skin products, such as calamine, to absorb the weeping of the skin caused by poisonous plants such as poison ivy, poison oak, and poison sumac. For relieving the pain or itching caused by these plants, the FDA document recommends a cold water compress and topical corticosteroids.
| Biology and health sciences | Specific drugs | Health |
2070483 | https://en.wikipedia.org/wiki/Spindle%20%28textiles%29 | Spindle (textiles) | A spindle is a straight spike, usually made from wood, used for spinning, twisting fibers such as wool, flax, hemp, cotton into yarn. It is often weighted at either the bottom, middle, or top, commonly by a disc or spherical object called a whorl; many spindles, however, are weighted simply by thickening their shape towards the bottom, e.g. Orenburg and French spindles. The spindle may also have a hook, groove, or notch at the top to guide the yarn. Spindles come in many different sizes and weights depending on the thickness of the yarn one desires to spin.
History
The origin of the first wooden spindle is lost to history because the materials did not survive. Whorl-weighted spindles date back at least to Neolithic times; spindle whorls have been found in archaeological digs around the world. Possible remains of spindle whorls were found in a Natufian village at Nahal Ein Gev II archeological site, Israel, from 12000 years ago.
A spindle is also part of traditional spinning wheels where it is horizontal, such as the Indian charkha and the great or walking wheel. In industrial yarn production, spindles are used as well; see spinning jenny, spinning mule and ring spinning.
The wood traditionally favoured for making spindles was that of Euonymus europaeus, from which derives the traditional English name spindle bush.
Hand spindles
Modern hand spindles fall into three basic categories: suspended spindles, supported spindles and grasped spindles. Supported and suspended spindles are normally held vertically, grasped spindles may be held vertically, horizontally or at an angle depending on the tradition.
Suspended spindles are so named because they are suspended to swing from the yarn after rotation has been started. Drop Spindles are a popular type of suspended spindle and get their name because the spindle is allowed to drop down while the thread is formed, allowing for a greater length of yarn to be spun before winding on. Suspended spindles also permit the spinner to move around while spinning, going about their day. However, there are practical limits to their size/weight.
Most supported spindles continue to rest with the tip on one's thigh, on the ground, on a table, or in a small bowl while rotating. Supported spindles come in a great variety of sizes, such as the very large, ~30" Navajo spindle, the small, extremely fast, metal takli for spinning cotton, and the tiniest Orenburg spindles (~20 cm, 15gm) for spinning gossamer lace yarns.
Grasped spindles are also known as: hand spindles, in the hand spindles, in hand spindles and twiddled spindles; there appears no consensus on nomenclature for this category of spindles though there have been various attempts at creating an agreed nomenclature including dividing this category of spindles into two, such as Crowfoot's attempts to define the difference between grasped and in hand spindles or merging this category into others, such as Franquemont's approach of classing them as supported spindles.
Grasped spindles remain held in the hand using finger or wrist movements to turn the spindle. French spindles are "twiddled" between the fingers of one hand while some types of Romanian spindles are grasped in the fist and turned through rotation of the wrist.
While spindle types are divided into these three main categories, some traditional spinning styles employ multiple or blended techniques. For example the Akha spindle, a short spindle with a large center-whorl disc, is supported by the hand of the spinner during drafting of cotton fibre, but during the adding of extra twist to stabilize the yarn, the spindle is dropped to rest on the yarn.
A familiar sight from history books is a spindle used in conjunction with a distaff, an upright stick with a large quantity of loose fibre wound around it, to be easily accessed. There are many other methods for controlling the pre-spun fibre, such as coiling it around one's lower arm, or through a bracelet, or wrapping it loosely around a yarn braid hanging from one's wrist.
Another way spindles are categorised is by the location of the whorl. The whorl, where present, may be located near the top, bottom or center of the spindle. For example a top-whorl drop spindle will have the whorl located near the top of the shaft underneath a hook that allows the spindle to suspend as it is being spun. The newly spun yarn is wound below the whorl and forms a ‘cop’. Depending on the location of the whorl and style of the spindle, the cop can be conical, football or ball shaped and it can be wound above, below or over the whorl.
Spindles can also be used for plying: intertwining two or more single strands of yarn together in order to create a stronger, more balanced, more durable yarn.
Anatomy
While hand spindles vary, there are some similarities in the parts that make up a spindle.
Shaft
Spindle shafts can be made out of a variety of materials such as wood, metal, bone or plastic. They may have very little shaping or be dramatically shaped enough to form part of the whorl. Shafts may be left plain or decorated with painting or carving.
The shaft is how the spinner inserts twist through turning it between the fingers or rolling it between the hand and another part of their anatomy, such as their thigh. The thickness of the shaft affects how fast the spindles spins, with narrower shafts causing a spindle to spin faster.
Many spindles will have a point at the top of the shaft to fix the thread to. Options include a simple length of shaft to tie the thread around, a shaped notch or bulb, or a hook.
Whorl
A whorl is a weight that is added to many types of spindles and can be made out of a large variety of materials including wood, metal, glass, plastic, stone, clay or bone. Whorls may be decorated or left plain, and they may be affixed permanently to the shaft or they may be removable.
Whorl shapes vary greatly and can include ball-shaped, disk-shaped and cross shaped whorls. The shape and mass distribution of the whorl affects the momentum it gives to the spindle while it is spinning. For example a center weighted whorl will spin very fast and short, while a rim-weighted disk-shaped whorl will spin longer and slower.
Whorls can be located near the top, middle or bottom of the spindle shaft. Whorl location can affect the stability of the spindle, with bottom whorl spindles being considered more stable.
Cop
The cop is not initially an intrinsic part of the spindle; however, as it is formed it plays a part in the spindle anatomy. Once a length of yarn or thread is spun it is wound around the spindle shaft or whorl to form a cop or a ball. As more yarn or thread is spun this makes the spindle heavier and more center-weighted, which has an effect on the momentum of the spindle. The overall shape of the cop and the skill in winding it also has an impact on how the spindle spins and how much thread or yarn can be stored on a spindle before it is "full".
Cops can be wound in a ball, cone or football shape.
Religious references
The spindle is closely associated with many goddesses, including the Germanic Holda, and Greek Artemis and Athena. It is often connected with fate, as the Greek Fates and the Norse Norns work with yarns that represent lives.
In Christianity, the apocryphal Gospel of James portrays the Virgin Mary as engaged in spinning thread for the Temple Curtain when the angel Gabriel tells her that she is going to bear Jesus.
Cultural references
Most modern illustrations of the fairy tale The Sleeping Beauty have the princess pricking her finger on the distaff of a modern flyer spinning wheel; this version of the wheel was either not invented or not as commonly used across Europe when the major versions of the story were commonly told. The Perrault and Brothers Grimm editions have her pricking her finger on a spindle, and some versions of the story have her pricking it on a great wheel, which is essentially a wheel powered spindle.
Biological references
In biology, the spindle apparatus is an assembly of proteins and DNA that forms during cell division to separate sister chromatids during mitosis or meiosis of eukaryotic cells. The word "mitosis" is derived from the Greek word "mitos," meaning warp thread.
Types
Supported spindles
Supported spindles spin on surfaces such as bowls that are usually ceramic or wooden. This type of spindle gives the user more control of the weight of the yarn. The various types of supported spindles range due to the difference in styles of spinning and yarn weight. Navajo spindles have longer shafts that should reach from the ground to the top of the thigh. The spun yarn is wound above the whorl. In Icelandic Viking times, the people used a high whorl lap spindle to spin wool into yarn. The high whorl lap spindle has a whorl on top with a hook on the top and a long shaft like other supported spindles. Unlike the Navajo spindle, however, the spun yarn is wrapped around the lower end of the spindle. The roving is prepped first by drafting it before the wool is spun into yarn. The separation between the drafting and spinning creates a softer yarn. Since it is a supported spindle, the yarn is less strained, which also makes the yarn soft.
Bottom whorl drop spindles
A typical bottom whorl spindle has the spun single ply yarn above the whorl instead of below it like top whorl drop spindles. Bottom whorl spindles are the preferred spindle for plying yarn. Turkish drop spindles have a low whorl spindle that have two arms that intersect to make a cross. Like other drop spindles it is spun and then it hangs while the wool is being drafted and the twist travels up the wool. When winding the spun yarn on to the arms it usually is put over and under the arms on the spindle. Even though it turns much more slowly than winding onto a regular spindle, it creates a center pull ball of yarn.
The Scottish drop spindle is called fairsaid, farsadh, or dealgan.
Top whorl drop spindles
Top whorl spindles commonly have a disc on the top of the shaft with a hook at the top.
| Technology | Spinning | null |
2071899 | https://en.wikipedia.org/wiki/Xiphactinus | Xiphactinus | Xiphactinus (from Latin and Greek for "sword-ray") is an extinct genus of large predatory marine ray-finned fish that lived during the late Albian to the late Maastrichtian. The genus grew up to in length, and superficially resembled a gargantuan, fanged tarpon. It is a member of the extinct order Ichthyodectiformes, which represent close relatives of modern teleosts.
The species Portheus molossus described by Cope is a junior synonym of X. audax. Skeletal remains of Xiphactinus have come from the Carlile Shale and Greenhorn Limestone of Kansas (where the first Xiphactinus fossil was discovered during the 1850s in the Niobrara Chalk), and Cretaceous formations all over the East Coast (most notably Georgia, Alabama, North Carolina, and New Jersey) in the United States, as well as Europe, Australia, the Kanguk and Ashville Formations of Canada, La Luna Formation of Venezuela and the Salamanca Formation in Argentina.
Paleobiology
Species of Xiphactinus were voracious predatory fish. At least a dozen specimens of X. audax have been collected with the remains of large, undigested or partially digested prey in their stomachs. In particular, one fossil "Fish-Within-A-Fish" specimen was collected by George F. Sternberg with another, nearly perfectly preserved long ichthyodectid Gillicus arcuatus inside of it. The larger fish apparently died soon after eating its prey, most likely owing to the smaller prey's struggling and rupturing an organ as it was being swallowed. This fossil is on display at the Sternberg Museum of Natural History in Hays, Kansas.
Like many other species in the Late Cretaceous oceans, a dead or injured individual was likely to be scavenged by sharks (Cretoxyrhina and Squalicorax). The remains of a Xiphactinus were found within a large specimen of Cretoxyrhina collected by Charles H. Sternberg. The specimen is on display at the University of Kansas Museum of Natural History.
Like modern tarpons, Xiphactinus likely spent its juvenile stage of life in shallow seaway margins for protection and to utilize rich food resources, possibly rare in open marine water, though this needs confirmation due to the lack of shallow, nearshore deposits from the Western Interior Seaway. The teeth of the juvenile specimen indicate that the diet of Xiphactinus probably didn't change notably during its growth, implying that even the small specimens would have been fish-eating predators.
"Unicerosaurus"
In 1982, a former Baptist minister, Carl Baugh, began excavations on the limestone beds of the Paluxy River, near Glen Rose, Texas, famous for its dinosaur tracks. Some of the tracks resembled human footprints and had been proclaimed since 1900 as evidence that dinosaurs and modern humans had once lived alongside one another. Scientists' investigations found the supposed human footprints to be "forms of elongate dinosaur tracks, while others were selectively highlighted erosional markings, and still others (on loose blocks) probable carvings." While excavating, he found a solitary "Y-shaped" fossil that he informally called "Unicerosaurus". In a 1987 popular article, John Armstrong described the fossil as a "Y-shaped petrified bone that appears to be the neural spine from a huge fish like the Portheus of Niobrara Chalk" that Baugh's museum "declared to be the forehead horn of a newly discovered dinosaur genus". The museum's exhibit told visitors that the "horn" belonged to "the unicorn of Job 38, one of three dinosaurs mentioned in Scripture; the others being behemoth and leviathan of Job 40 and 41", and that the horn was able to fold back like the blade of a jack knife. Although some Young Earth Creationists shared Baugh's interpretations of the biblical Behemoth and Leviathan, Baugh's claims were not taken seriously either by Christian organizations or the scientific community.
In popular culture
In October 2010, Kansas House Rep. Tom Sloan (R-Lawrence) announced that he would introduce legislation to make Xiphactinus audax, a.k.a. the "X-fish", the state fossil of Kansas. Ultimately, Tylosaurus was selected instead.
| Biology and health sciences | Prehistoric osteichthyans | Animals |
21476681 | https://en.wikipedia.org/wiki/Scimitar | Scimitar | A scimitar ( or ) is a single-edged sword with a convex curved blade of about 76.2 to 91.44cm (30 to 36 inches) associated with Middle Eastern, South Asian, or North African cultures. A European term, scimitar does not refer to one specific sword type, but an assortment of different Eastern curved swords inspired by types introduced to the Middle East by Central Asian ghilmans. These swords include the Persian shamshir (the origin of the word scimitar), the Arab saif, the Indian talwar, the North African nimcha, the Turkish kilij, and the Afghan pulwar. All such swords are originally derived from earlier curved swords developed in Turkic Central Asia (Turkestan).
Etymology
The English term scimitar is attested from the mid-16th century and derives partly from the Middle French cimeterre (15th century) and partly the Italian scimitarra. The ultimate source of these terms is corruptions of the Persian shamshir. Scimitar became used to describe all curved oriental blades, in contrast to the straight and double-edged European swords of the time.
The term saif in Arabic can refer to any Middle Eastern (or North African, South Asian) sword, straight or curved. Saif cognates with the ancient Greek xiphos, which may have been borrowed from a Semitic language, as both saif and xiphos go back to an old (Bronze Age) wanderwort of the Eastern Mediterranean of unknown ultimate origin. Richard F. Burton derives both words from the Egyptian sfet.
History of use
The earliest evidence of scimitars is from the 9th century among soldiers in Khurasan. They were used in horse warfare because of their relatively light weight when compared to larger swords and their curved design, good for slashing opponents while riding on a horse. Nomadic horsemen learned from experience that a curved edge is better for cutting strikes because the arc of the blade matches that of the sweep of the rider's arm as they slash the target while galloping. Turks, Mongols, Rajputs and Sikhs used scimitars in warfare, among many other peoples.
The scimitar was widespread throughout the Middle East from at least the Ottoman period until the age of smokeless powder firearms relegated swords to dress and ceremonial function. The Egyptian khopesh, brought to Egypt by the Hyksos, resembled scimitars. The khopesh is sometimes considered a scimitar. Early swords in Islamic lands were typically straight and double-edged, following the tradition of the weapons used by the Islamic prophet Muhammad. Though the famous double-edged sword, Zulfiqar wielded by Ali was of a curved design, the curved design was probably introduced into central Islamic lands by Turkic warriors from central Asia who were employed as royal body-guards in the 9th century and an Abbasid era blade has been discovered from Khurasan. These Turkic warriors sported an early type of sabre which had been used in central Asia since the 7th century, but failed to gain wider appeal initially in Islamic lands. There is a single surviving Seljuk saber from approximately the year 1200, which may indicate that under that empire curved blades saw some popularity. Following the Mongol invasions of the 13th century the curved swords favored by the Turkic cavalry, formed lasting impacts across much of the Middle East. The adoption of these swords was incremental, starting not long after Mongol conquest, and lasting well into the 15th century. During Islamization of the Turks, the kilij became more and more popular in the Islamic armies. When the Seljuk Empire invaded Persia and became the first Turkic Muslim political power in Western Asia, kilij became the dominant sword form. The Iranian shamshir was created during the Turkic Seljuk Empire period of Iran.
Symbolism
The sword (or saif) is an important symbol in Arab cultures, and is used as a metaphor in many phrases in the Arabic language. The word occurs also in various symbolic and status titles in Arabic (and adopted in other languages) used in Islamic states, notably:
In the Yemenite independent imamate:
Saif al-Haqq, meaning "Sword of Truth".
Saif al-Islam "Sword of submission to Allah" or "Sword of Islam", was a subsidiary title borne (after their name and patronym) by male members of the al-Qasimi dynasty (whose primary title, before the name, was Amir), especially sons of the ruling Imam.
Sayf al-Dawla and variations mean "Sword of the State".
Saif Ullah Al-masloul the "drawn sword of God" was conferred by Muhammad, uniquely, to the recent convert and military commander Khalid ibn al-Walid.
Saif ul-Mulk "sword of the realm" was an honorary title awarded by the Mughal Padshahs of Hind (India), e.g. as one of the personal titles (including Nawab bahadur, one rank above his dynasty's) conferred in 1658 by the Mughal emperor Aurangzeb to Nawab Muhammad Bayazid Khan Bahadur, a high mansabdar, whose jagir of Malerkotla was by sanad raised to Imperial riyasat, thus becoming an independent ruler.
Saif ul-Ali, "Sword of Ali", referring to arguably most famous sword in Islamic history, belonging to both Muhammad, and later, Ali, Zulfikar, and with which Ali slew a Makkan foot soldier, cleaving both his helmet and head, at the Battle of Uhud, and with which he (Ali) slew Amr, a ferocious and devastating Makkan soldier at the Battle of the Trench at Madinah.
Saif and Saif al Din "Sword of the religion" are also common masculine (and male) Islamic names.
The scimitar appears as a symbol of the Russians in the Finnish coat of arms of the Province of Karelia, which depicts two armored arms fighting with swords – one Western and one Eastern, representing Karelia's troubled history as a border region between the east and the west. From this context, the sword and scimitar have found their way into the coat of arms of Finland, which depicts a lion brandishing a sword and trampling a scimitar.
In Islamophobia: Making Muslims the Enemy, Peter Gottschalk and Gabriel Greenberg argue that the scimitar has been appropriated by Western culture and Hollywood to symbolize Arab Muslims in a negative light. Even though Muslims used straight-edge swords for the first two centuries of the Crusades, European Christians may have more closely tied the Christian cross-like shape of the swords to their cause. The authors commented that American cartoonists use the scimitar to symbolize "Muslim barbarity", despite the irony in scimitars being worn with some American military uniforms.
In European art
In Shakespeare's works, the scimitar was a symbol for the East and the Islamic world.
Scimitars were used in 19th century orientalist depictions of Middle Eastern men. In the 20th century, they were often used to indicate that a character was Middle Eastern and occupying a villain role.
In media
Scimitars were used by Calormene warriors and royalty in The Chronicles of Narnia, such as The Horse and His Boy and The Last Battle.
In the 1975 film Monty Python and the Holy Grail, a character says: "If I went 'round, sayin' I was an emperor, just because some moistened bint had lobbed a scimitar at me, they'd put me away!" This is in satirical reference to the legend of King Arthur and the Lady of the Lake.
A scimitar was used by Morgan Freeman when playing Azeem Edin Bashir Al-Bakir, the Moorish ally to the title character of Robin Hood: Prince of Thieves.
Scimitars were weapons used by Aladdin in the 1992 film.
Scimitars are used extensively in the Assassin's Creed video game series, particularly in the first game, Assassin's Creed II, Assassin's Creed Brotherhood, Assassin's Creed: Revelations, Assassin's Creed III: Liberation, Assassin's Creed IV: Black Flag, Assassin's Creed Rogue and Assassin's Creed Unity.
Scimitars are a popular weapon in the RuneScape franchise.
| Technology | Swords | null |
21481546 | https://en.wikipedia.org/wiki/Wheelchair | Wheelchair | A wheelchair is a mobilized form of chair using 2 or more wheels, a footrest, and an armrest usually cushioned. It is used when walking is difficult or impossible to do due to illnesses, injury, disabilities, or age-related health conditions. Wheelchairs provide mobility, postural support, and freedom to those who cannot walk or have difficulty walking, enabling them to move around, participate in everyday activities, and live life on their own terms. []
Wheelchairs come in a wide variety of formats to meet the specific needs of their users. They may include specialized seating adaptions, and individualized controls, and may be specific to particular activities, as with sports wheelchairs and beach wheelchairs. The most widely recognized distinction is between motorized wheelchairs, where propulsion is provided by batteries and electric motors, and manual wheelchairs, where the propulsive force is provided either by the wheelchair user or occupant pushing the wheelchair by hand (self-propelled), by an attendant pushing from the rear using the handle(s), or by an attendant pushing from the side use a handle attachment.
History
The Sanskrit grammarian Pāṇini (circa 5th century BCE India) mentions the Sanskrit word parpa for wheelchair in his Aṣṭādhyāyī 4.4.10.
The earliest records of wheeled furniture are an inscription found on a stone slate in China and a child's bed depicted in a frieze on a Greek vase, both dating between the 6th and 5th century BC. The first records of wheeled seats being used for transporting disabled people date to three centuries later in China; the Chinese used early wheelbarrows to move people as well as heavy objects. A distinction between the two functions was not made for another several hundred years, until around AD 525, when images of wheeled chairs made specifically to carry people began to occur in Chinese art.
Although Europeans eventually developed a similar design, this method of transportation did not exist until 1595 when an unknown inventor from Spain built one for King Phillip II. Although it was an elaborate chair having both armrests and leg rests, the design still had shortcomings since it did not feature an efficient propulsion mechanism and thus required assistance to propel it. This makes the design more comparable to a modern-day highchair or portable throne for the wealthy than to a modern-day wheelchair for disabled people.
In 1655, Stephan Farffler, a 22-year-old paraplegic watchmaker, built the world's first self-propelling chair on a three-wheel chassis using a system of cranks and cogwheels. However, the device resembled a hand bike more than a wheelchair since the design included hand cranks mounted at the front wheel.
A self-propelled wheelchair was made for the Parliamentarian commander-in-chief Sir Thomas Fairfax due to the many injuries he had received during the English Civil War, and he used it during the final years of his life. The wheelchair of Thomas Fairfax is currently on display at the National Civil War Centre in Newark-on-Trent.
The invalid carriage or Bath chair brought the technology into more common use from around 1760.
In 1887, wheelchairs ("rolling chairs") were introduced to Atlantic City so invalid tourists could rent them to enjoy the Boardwalk. Soon, many healthy tourists also rented the decorated "rolling chairs" and servants to push them as a show of decadence and treatment they could never experience at home.
Modern wheelchairs
In 1933 Harry C. Jennings Sr. and his disabled friend Herbert Everest, both mechanical engineers, invented the first lightweight, steel, folding, portable wheelchair. Everest had previously broken his back in a mining accident. Everest and Jennings saw the business potential of the invention and went on to become the first mass-market manufacturer of wheelchairs. Their "X-brace" design is still in common use, albeit with updated materials and other improvements. The X-brace idea came to Jennings from the men's folding "camp chairs/stools", rotated 90 degrees, used in the outdoors and at the mines.
Types
There are a wide variety of types of wheelchairs, differing by propulsion method, mechanisms of control, and technology used. Some wheelchairs are designed for general everyday use, others for single activities, or to address specific access needs. Innovation within the wheelchair industry is relatively common, but many innovations ultimately fall by the wayside, either from over-specialization or from failing to come to market at an accessible price point. The iBOT is perhaps the best-known example of this in recent years.
Manual self-propelled wheelchairs
A self-propelled manual wheelchair incorporates a frame, seat, one or two footplates (footrests), and four wheels: usually two caster wheels at the front and two large wheels at the back. There will generally also be a separate seat cushion. The larger rear wheels usually have push-rims of slightly smaller diameter projecting just beyond the tyre; these allow the user to manoeuvre the chair by pushing on them without requiring them to grasp the tyres. Manual wheelchairs generally have brakes that bear on the tyres of the rear wheels, however, these are solely a parking brakes and in-motion braking is provided by the user's palms bearing directly on the push-rims. As this causes friction and heat build-up, particularly on long downslopes, many wheelchair users will choose to wear padded wheelchair gloves. Manual wheelchairs often have two push handles at the upper rear of the frame to allow for manual propulsion by a second person, however many active wheelchair users will remove these to prevent unwanted pushing from people who believe they are being helpful.
Everyday manual wheelchairs come in two major varieties, folding or rigid. Folding chairs are generally low-end designs, whose predominant advantage is being able to fold, generally by bringing the two sides together. This is an advantage for people who need to store the wheelchair frequently or to put it in a small vehicle. Rigid wheelchairs have permanently welded joints and many fewer moving parts. This reduces the energy required to push the chair by eliminating many points where the chair would flex and absorb energy as it moves. Welded rather than folding joints also reduce the overall weight of the chair. Rigid chairs typically feature instant-release rear wheels and backrests that fold down flat, allowing the user to dismantle the chair quickly for storage in a car. A few wheelchairs attempt to combine the features of both designs by providing a fold-to-rigid mechanism in which the joints are mechanically locked when the wheelchair is in use.
Many rigid models are made with light materials such as aluminium and titanium, and wheelchairs of composite materials such as carbon-fibre have started to appear. Ultra lightweight rigid wheelchairs are commonly known as 'active user chairs' as they are ideally suited to independent use. Another innovation in rigid chair design is the installation of shock absorbers, such as "Frog Legs", which cushion the bumps over which the chair rolls. These shock absorbers may be added to the front wheels, to the rear wheels, or both. Rigid chairs also have the option for their rear wheels to have a camber, or tilt, which angles the tops of the wheels in toward the chair. This allows for more mechanically efficient propulsion by the user and also makes it easier to hold a straight line while moving across a slope. Sport wheelchairs often have large camber angles to improve stability.
Rigid-framed chairs are generally made to measure, to suit both the specific size of the user and their needs and preferences around areas such as the "tippyness" of the chair - determined by the distance between the center of gravity and the rear axle. Experienced users with sufficient upper-body strength can generally balance the chair on its rear wheels, a "wheelie", and the "tippyness" of the chair controls the ease with which this can be initiated. The wheelie allows an independent wheelchair user to climb and descend curbs and move more easily over small obstacles and irregular ground such as cobbles.
The rear wheels of self-propelled wheelchairs typically range from in diameter, and commonly resemble bicycle wheels. Wheels are rubber-tired and may be solid, pneumatic or gel-filled. The wheels of folding chairs may be permanently attached, but those for rigid chairs are commonly fitted with quick-release axles activated by depressing a button at the centre of the wheel.
All major varieties of wheelchairs can be highly customized for the user's needs. Such customization may encompass the seat dimensions, height, seat angle, footplates, leg rests, front caster outriggers, adjustable backrests and controls. Various optional accessories are available, such as anti-tip bars or wheels, safety belts, adjustable backrests, tilt and/or recline features, extra support for limbs or head and neck, holders for crutches, walkers or oxygen tanks, drink holders, and mud and wheel-guards as clothing protectors.
Light weight and high costs are related to the manual wheelchair market. At the low-cost end, heavy, folding steel chairs with sling seats and little adaptability dominate. Users may be temporarily disabled, or using such a chair as a loaner, or simply unable to afford better. These chairs are common as "loaners" at large facilities such as airports, amusement parks and shopping centers. A slightly higher price band sees the same folding design produced in aluminium. The high end of the market contains ultra-light models, extensive seating options and accessories, all-terrain features, and so forth. The most expensive manual chairs may rival the cost of a small car.
Manual attendant-propelled wheelchairs
An attendant-propelled wheelchair (also known as a companion or transfer chair) is generally similar to a self-propelled manual wheelchair, but with small diameter wheels at both front and rear. The chair is maneuvered and controlled by a person standing at the rear and pushing on handles incorporated into the frame. Braking is supplied directly by the attendant who will usually also be provided with a foot- or hand-operated parking brake.
These chairs are common in institutional settings and as loaner-chairs in large public venues. They are usually constructed from steel as light weight is less of a concern when the user is not required to self-propel.
Specially designed transfer chairs are now required features at airports in much of the developed world in order to allow access down narrow airliner aisles and facilitate the transfer of wheelchair-using passengers to and from their seats on the aircraft.
Powered wheelchairs
An electric-powered wheelchair, commonly called a "powerchair" is a wheelchair that additionally incorporates batteries and electric motors into the frame and that is controlled by either the user or an attendant, most commonly via a small joystick mounted on the armrest, or on the upper rear of the frame. Alternatives exist for the traditional manual joystick, including head switches, chin-operated joysticks, sip-and-puff controllers or other specialist controls, which may allow independent operation of the wheelchair for a wider population of users with varying motor impairments. Ranges of over 10 miles/15 km are commonly available from standard batteries.
Powerchairs are commonly divided by their access capabilities. An indoor-chair may only reliably be able to cross completely flat surfaces, limiting them to household use. An indoor-outdoor chair is less limited, but may have restricted range or ability to deal with slopes or uneven surfaces. An outdoor chair is more capable, but will still have a very restricted ability to deal with rough terrain. A very few specialist designs offer a true cross-country capability.
Powerchairs have access to the full range of wheelchair options, including ones that are difficult to provide in an unpowered manual chair, but have the disadvantage of significant extra weight. Where an ultra-lightweight manual chair may weigh under 10 kg, the largest outdoor power-chairs may weigh 200 kg or more.
Smaller power chairs often have four wheels, with front or rear wheel drive, but large outdoor designs commonly have six wheels, with small wheels at front and rear and somewhat larger powered wheels in the centre.
A power-assisted wheelchair is a recent development that uses the frame and seating of a typical rigid manual chair while replacing the standard rear wheels with wheels of similar size which incorporate batteries and battery-powered motors in the hubs. A floating rim design senses the pressure applied by the user's push and activates the motors proportionately to provide a power assist. This results in the convenience, and small size of a manual chair while providing motorised assistance for rough/uneven terrain and steep slopes that would otherwise be difficult or impossible to navigate, especially by those with limited upper-body function. As the wheels necessarily come at a weight penalty it is often possible to exchange them with standard wheels to match the capabilities of the wheelchair to the current activity.
Mobility scooters
Mobility scooters share some features with powerchairs, but primarily address a different market segment, people with a limited ability to walk, but who might not otherwise consider themselves disabled. Smaller mobility scooters are typically three wheeled, with a base on which is mounted a basic seat at the rear, with a control tiller at the front. Larger scooters are frequently four-wheeled, with a much more substantial seat.
Opinions are often polarized as to whether mobility scooters should be considered wheelchairs or not, and negative stereotyping of scooter users can be worse than for some manual or power-chair users. Some commercial organisations draw a distinction between power-chairs and scooters when making access provisions due to a lack of clarity in the law as to whether scooters fall under the same equality legislation as wheelchairs.
Single-arm drive wheelchairs
One-arm or single arm drive enables a user to self-propel a manual wheelchair using only a single arm. The large wheel on the same side as the arm to be used is fitted with two concentric handrims, one of smaller diameter than the other. On most models the outer, smaller rim, is connected to the wheel on the opposite side by an inner concentric axle. When both handrims are grasped together, the chair may be propelled forward or backward in a straight line. When either handrim is moved independently, only a single wheel is used and the chair will turn left or right in response to the handrim used. Some wheelchairs, designed for use by hemiplegics, provide a similar function by linking both wheels rigidly together and using one of the footplates to control steering via a linkage to the front caster.
Reclining and tilting wheelchairs
Reclining or tilt-in-space wheelchairs have seating surfaces that can be tilted to various angles. The original concept was developed by an orthotist, Hugh Barclay, who worked with disabled children and observed that postural deformities such as scoliosis could be supported or partially corrected by allowing the wheelchair user to relax in a tilted position. The feature is also of value to users who are unable to sit upright for extended periods for pain or other reasons.
In the case of reclining wheelchairs, the seat-back tilts back, and the leg rests can be raised, while the seat base remains in the same position, somewhat similar to a common recliner chair. Some reclining wheelchairs lean back far enough that the user can lie down completely flat. Reclining wheelchairs are preferred in some cases for some medical purposes, such as reducing the risk of pressure sores, providing passive movement of hip and knee joints, and making it easier to perform some nursing procedures, such as intermittent catheterization to empty the bladder and transfers to beds, and also for personal reasons, such as people who like using an attached tray. The use of reclining wheelchairs is particularly common among people with spinal cord injuries such as quadriplegia.
In the case of tilting wheelchairs, the seat-back, seat base, and leg rests tilt back as one unit, somewhat similar to the way a person might tip a four-legged chair backward to balance it on the back legs. While fully reclining spreads the person's weight over the entire back side of the body, tilting wheelchairs transfer it from only the buttocks and thighs (in the seated position) to partially on the back and head (in the tilted position). Tilting wheelchairs are preferred for people who use molded or contoured seats, who need to maintain a particular posture, who adversely affected by sheer forces (reclining causes the body to slide slightly every time), or who need to keep a communication device, powered wheelchair controls, or other attached device in the same relative position throughout the day. Tilting wheelchairs are commonly used by people with cerebral palsy, people with some muscle diseases, and people with limited range of motion in the hip or knee joints. Tilting options are more common than reclining options in wheelchairs designed for use by children.
Standing wheelchairs
A standing wheelchair is one that supports the user in a nearly standing position. They can be used as both a wheelchair and a standing frame, allowing the user to sit or stand in the wheelchair as they wish. Some versions are entirely manual, others have a powered stand on an otherwise manual chair, while others have full power, tilt, recline and variations of powered stand functions available. The benefits of such a device include, but are not limited to: aiding independence and productivity, raising self-esteem and psychological well-being, heightening social status, extending access, relief of pressure, reduction of pressure sores, improved functional reach, improved respiration, reduced occurrence of UTI, improved flexibility, help in maintaining bone mineral density, improved passive range motion, reduction in abnormal muscle tone and spasticity, and skeletal deformities. Other wheelchairs provide some of the same benefits by raising the entire seat to lift the user to standing height.
Sports wheelchairs
A range of disabled sports have been developed for disabled athletes, including basketball, rugby, tennis, racing and dancing. The wheelchairs used for each sport have evolved to suit the specific needs of that sport and often no longer resemble their everyday cousins. They are usually non-folding (in order to increase rigidity), with a pronounced negative camber for the wheels (which provides stability and is helpful for making sharp turns), and often are made of composite, lightweight materials. Even seating positions may be radically different, with racing wheelchairs generally used in a kneeling position. Sport wheelchairs are rarely suited for everyday use, and are often a 'second' chair specifically for sports use, although some users prefer the sports options for everyday use. Some disabled people, specifically lower-limb amputees, may use a wheelchair for sports, but not for everyday activities.
Motorized wheelchair sports
While most wheelchair sports use manual chairs, some power chair sports, such as powerchair football, exist. Hockey can also be played from electrical wheelchairs.
Wheelchair stretchers
Wheelchair stretchers are a variant of wheeled stretchers/gurneys that can accommodate a sitting patient, or be adjusted to lie flat to help in the lateral (or supine) transfer of a patient from a bed to the chair or back. Once transferred, the stretcher can be adjusted to allow the patient to assume a sitting position.
All-terrain wheelchairs
All-terrain wheelchairs can allow users to access terrain otherwise completely inaccessible to a wheelchair user. Two different formats have been developed. One hybridises wheelchair and mountain bike technology, generally taking the form of a frame within which the user sits and with four mountain bike wheels at the corners. In general, there are no push-rims and propulsion/braking is by pushing directly on the tyres.
A more common variant is the beach wheelchair (beach-going wheelchair) which can allow better mobility on beach sand, including in the water, on uneven terrain, and even on snow. The common adaptation among the different designs is that they have extra-wide balloon wheels or tires, to increase stability and decrease ground pressure on uneven or unsteady terrain. Different models are available, both manual and battery-driven. In some countries in Europe, where accessible tourism is well established, many beaches have wheelchairs of this type available for loan/hire.
Smart wheelchairs
A smart wheelchair is any powerchair using a control system to augment or replace user control. Its purpose is to reduce or eliminate the user's task of driving a powerchair. Usually, a smart wheelchair is controlled via a computer, has a suite of sensors and applies techniques in mobile robotics, but this is not necessary. The type of sensors most frequently used by smart wheelchairs are the ultrasonic acoustic range finder (i.e. sonar) and infrared red (IR) range finder. The interface may consist of a conventional wheelchair joystick, a "sip-and-puff" device or a touch-sensitive display. This differs from a conventional powerchair, in which the user exerts manual control over speed and direction without intervention by the wheelchair's control system.
Smart wheelchairs are designed for a variety of user types. Some are designed for users with cognitive impairments, such as dementia, these typically apply collision-avoidance techniques to ensure that users do not accidentally select a drive command that results in a collision. Others focus on users living with severe motor disabilities, such as cerebral palsy, or with quadriplegia, and the role of the smart wheelchair is to interpret small muscular activations as high-level commands and execute them. Such wheelchairs typically employ techniques from artificial intelligence, such as path-planning.
Technological developments
Recent technological advances are slowly improving wheelchair and powerchair technology.
A variation on the manually-propelled wheelchair is the Leveraged Freedom Chair (LFC), designed by the MIT Mobility Lab. This wheelchair is designed to be low-cost, constructed with local materials, for users in developing countries. Engineering modifications have added hand-controlled levers to the LFC, to enable users to move the chair over uneven ground and minor obstacles, such as bumpy dirt roads, that are common in developing countries. It is under development, and has been tested in Kenya and India so far.
The addition of geared, all-mechanical wheels for manual wheelchairs is a new development incorporating a hypocycloidal reduction gear into the wheel design. The 2-gear wheels can be added to a manual wheelchair. The geared wheels provide a user with additional assistance by providing leverage through gearing (like a bicycle, not a motor). The two-gear wheels offer two speed ratios- 1:1 (no help, no extra torque) and 2:1, providing 100% more hill climbing force. The low gear incorporates an automatic "hill hold" function which holds the wheelchair in place on a hill between pushes, but will allow the user to override the hill hold to roll the wheels backward if needed. The low gear also provides downhill control when descending.
A recent development related to wheelchairs is the handcycle. They come in a variety of forms, from the road and track racing models to off-road types modelled after mountain bikes. While dedicated handcycle designs are manufactured, clip-on versions are available that can convert a manual wheelchair to a handcycle in seconds. The general concept is a clip-on front-fork with hand-pedals, usually attaching to a mounting on the footplate. A somewhat related concept is the Freewheel, a large dolley wheel attaching to the front of a manual wheelchair, again generally to the footplate mounting, which improves wheelchair performance over rough terrain. Unlike a handcycle, a wheelchair with a Freewheel continues to be propelled via the rear wheels. There are several types of hybrid-powered handcycles where hand-pedals and used along with the electrical motor that helps on hills and large distances.
The most recent generation of clip-on handcycles is fully electrical wheelchair power add-ons that use lithium-ion battery, brushless DC electric motor and light-weight aluminium frames with easy to attach clamps to convert almost any manual wheelchair into electrical trike in seconds. That makes long-distance journeys and everyday tasks much easier and keeps wheelchair users hands clean.
There have been significant efforts over the past 20 years to develop stationary wheelchair trainer platforms that could enable wheelchair users to exercise as one would on a treadmill or bicycle trainer.
Some devices have been created that could be used in conjunction with virtual travel and interactive gaming similar to an omnidirectional treadmill. This convergence of virtual reality and a treadmill have been used for pediatric and adult rehabilitation to regain walking skills.
In 2011, British inventor Andrew Slorance developed Carbon Black the first wheelchair to be made almost entirely out of carbon fiber
Recently, EPFL's CNBI project has succeeded in making wheelchairs that can be controlled by brain impulses.
Interest in electric-powered wheelchairs that are able to climb stairs has increased over the past twenty years. Therefore, many electric wheelchairs with the ability to climb stairs have been developed. Electric-powered wheelchairs with climbing ability need to be stronger and have greater movement in comparison to an electric-powered wheelchair that cannot climb stairs. They must also be stable in order to prevent injury to the wheelchair user. There are currently a number of electric powered wheelchairs that are able to climb stairs available to purchase. Technical developments are continuing in this area.
Experiments have also been made with unusual variant wheels, like the omniwheel or the Mecanum wheel. These allow for a broader spectrum of movement, but have made no mass-market penetration.
The electric wheelchair shown on the right is fitted with Mecanum wheels (sometimes known as Ilon wheels) which give it complete freedom of movement. It can be driven forwards, backward, sideways, and diagonally, and also turned around on the spot or turned around while moving, all operated from a simple joystick.
Other variants
Foot propulsion of a manual wheelchair by the occupant is possible for users who have limited hand movement capabilities or simply do not wish to use their hands for propulsion. Foot propulsion also allows patients to exercise their legs to increase blood flow and limit further disability. Users who do this commonly may elect to have a lower seat height and no footplate to better suit the wheelchair to their needs.
Wheelbase chairs are powered or manual wheelchairs with especially molded seating systems interfaced with them for users with a more complicated posture. A molded seating system involves taking a cast of a person's best achievable seated position and then either carving the shape from memory foam or forming a plastic mesh around it. This seat is then covered, framed, and attached to a wheelbase.
A bariatric wheelchair is one designed to support larger weights; most standard chairs are designed to support no more than 250 lb (113 kg) on average.
Pediatric wheelchairs are another available subset of wheelchairs. These can address needs such as being able to play on the floor with other children, or cater for children in large hip-spica casts due to problems such as hip dysplasia.
Hemi wheelchairs have lower seats which are designed for easy foot propulsion. The decreased seat height also allows them to be used by children and shorter individuals.
A knee scooter is a related device with some features of a wheelchair and some walking aids. Unlike wheelchairs they are only suitable for below knee injuries to a single leg. The user rests the injured leg on the scooter, grasps the handlebars, and pushes with the uninjured leg.
Some walkers can be used as a wheelchair. These walkers have seat and foot plates, so an attendant can push while the patient is sitting on the walker. This is useful for a person who gets tired while walking with a walker, or has a limited walking range meaning the person can walk, but after a while, the person will collapse and fall to the ground.
A commode wheelchair is a wheelchair made for the bathroom. A commode wheelchair has a hole in the seat so the user does not have to transfer into the toilet. Sometimes the hole can be covered. Sometimes there is a pan attached to the hole, so the user can urinate/defecate without having to wheel over the toilet.
Mobility and access
Buildings
Adapting the built environment to make it more accessible to wheelchair users is one of the key campaigns of disability rights movements and local equality legislation such the Americans with Disabilities Act of 1990 (ADA). The social model of disability defines 'disability' as the discrimination experienced by people with impairments as a result of the failure of society to provide the adaptions needed for them to participate in society as equals. This includes both physical adaptions of the built environment and adaption of organizational and social structures and attitudes. A core principle of access is universal design - that all people regardless of disability are entitled to equal access to all parts of society like public transportation and buildings. A wheelchair user is less disabled in an environment without stairs.
Access starts outside of the building, with the provision of reduced height kerb-cuts where wheelchair users may need to cross roads, and the provision of adequate wheelchair parking, which must provide extra space in order to allow wheelchair users to transfer directly from seat to chair. Some tension exists between access provisions for visually impaired pedestrians and wheelchair users and other mobility impaired pedestrians as textured paving, vital for visually impaired people to recognise the edge of features such as light-controlled crossings, is uncomfortable at best, and dangerous at worst, to those with mobility impairments.
For access to public buildings, it is frequently necessary to adapt older buildings with features such as ramps or elevators in order to allow access by wheelchair users and other people with mobility impairments. Other important adaptations can include powered doors, lowered fixtures such as sinks and water fountains, and accessible toilets with adequate space and grab bars to allow the disabled person to transfer out of their wheelchair onto the fixture. Access needs for people with other disabilities, for instance visual impairments, may also be required, such as by provision of high visibility markings on the edges of steps and braille labelling. Increasingly new construction for public use is required by local equality laws to have these features incorporated at the design stage.
The same principles of access that apply to public buildings also apply to private homes and may be required as part of local building regulations. Important adaptations include external access, providing sufficient space for a wheelchair user to move around the home, doorways that are wide enough for convenient use, access to upper floors, where they exist, which can be provided either by dedicated wheelchair lifts, or in some cases by using a stairlift to transfer between wheelchairs on different floors, and by providing accessible bathrooms with showers and/or bathtubs that are designed for accessibility. Accessible bathrooms can permit the use of mobile shower chairs or transfer benches to facilitate bathing for people with disabilities. Wet rooms are bathrooms where the shower floor and bathroom floor are one continuous waterproof surface. Such floor designs allow a wheelchair user using a dedicated shower chair, or transferring onto a shower seat, to enter the shower without needing to overcome a barrier or lip.
The construction of low floor trams and buses are increasingly required by law, whereas the use of inaccessible features such as paternoster lifts in public buildings without any alternative methods of wheelchair access is increasingly deprecated. Modern architecture is increasingly required by law and recognised good practise to incorporate better accessibility at the design stage.
In many countries, such as the UK, the owners of inaccessible buildings who have not provided permanent access measures are still required by local equality legislation to provide 'reasonable adjustments' to ensure that disabled people are able to access their services and are not excluded. These may range from keeping a portable ramp on hand to allow a wheelchair user to cross an inaccessible threshold, to providing personal service to access goods they are not otherwise able to reach.
Vehicles
Public transit vehicles are increasingly required to be accessible to people who use wheelchairs.
In the UK, all single deck buses are required to be accessible to wheelchair users by 2017, all double-deck coaches by 2020. Similar requirements exist for trains, with most trains already incorporating a number of wheelchair-spaces.
The EU has required airline and airport operators to support the use of airports and airliners by wheelchair users and other 'Persons with Reduced Mobility' since the introduction of EU Directive EC1107/2006.
In Los Angeles there is a program to remove a small amount of seating on some trains to make more room for bicycles and wheelchairs.
New York City's entire bus system is wheelchair-accessible, and a multimillion-dollar renovation program is underway to provide elevator access to many of the city's 485 subway stations.
In Adelaide, Australia, all public transport has provision for at least two wheelchairs per bus, tram or train. In addition, all trains have space available for bicycles.
The Washington, D.C. Metro system features complete accessibility on all its subways and buses.
In Paris, France, the entire bus network, i.e. 60 lines, has been accessible to wheelchair users since 2010.
In the United States, a wheelchair that has been designed and tested for use as a seat in motor vehicles is often referred to as a "WC19 Wheelchair" or a "transit wheelchair". ANSI-RESNA WC19 (officially, SECTION 19 ANSI/RESNA WC/VOL. 1 Wheelchairs for use in Motor Vehicles) is a voluntary standard for wheelchairs designed for use when traveling facing forward in a motor vehicle. ISO 7176/19 is an international transit wheelchair standard that specifies similar design and performance requirements as ANSI/RESNA WC19.
There are special vans equipped for wheelchairs. These vans are large and have a ramp on a side door or the back door, so a wheelchair can get inside the vehicle while the user is still in it. Some of the back seats will be removed and replaced with wheelchair security harnesses. Sometimes wheelchair vans are equipped so the wheelchair user can drive the van without getting out of the wheelchair.
A vehicle can be equipped with hand controls. Hand controls are used when a person can not move their legs to push the pedals. The hand controls do the pushing of the pedals. Some racecar drivers are paralyzed and use hand controls.
Public realm
Some elements of the outdoor public realm present barriers to wheelchair use, such as grade-separated areas where no ramps are provided. UK guidance recommends a minimum width of 1.5m for a wheelchair user and an ambulant pedestrian to move side by side; with maximum ramp gradients of 8% for manual wheelchair use (5% is preferred).
Distribution organizations
Several organizations exist that help to give and receive wheelchair equipment. Organizations that accept wheelchair equipment donations typically attempt to identify recipients and match them with the donated equipment they have received. Organizations that accept donations in the form of money for wheelchairs typically have the wheelchairs manufactured and distributed in large numbers, often in developing countries. Organizations focusing on wheelchairs include Direct Relief, the Free Wheelchair Mission, Hope Haven, Personal Energy Transportation, the Wheelchair Foundation and WheelPower.
In the United Kingdom wheelchairs are supplied and maintained free of charge for non-ambulatory disabled people.
Seating systems
Wheelchair seating systems are designed both to support the user in the sitting position and to redistribute pressure from areas of the body that are at risk of pressure ulcers. For someone in the sitting position, the parts of the body that are the most at risk for tissue breakdown include the ischial tuberosities, coccyx, sacrum and greater trochanters. Wheelchair cushions are the prime method of delivering this protection and are nearly universally used. Wheelchair cushions are also used to provide stability, comfort, aid posture and absorb shock. Wheelchair cushions range from simple blocks of foam costing a few pounds or dollars, to specifically engineered multilayer designs with costs running into the hundreds of dollars.
Prior to 1970, little was known about the effectiveness of wheelchair cushions and there was not a clinical method of evaluating wheelchair seat cushions. Most recently, pressure imaging (or pressure mapping) is used to help determine each individual's pressure distribution to properly determine and fit a seating system.
While almost all wheelchair users will use a wheelchair cushion, some users need more extensive postural support. This can be provided by adaptions to the back of the wheelchair, which can provide increased rigidity, head/neck rests and lateral support and in some cases by adaptions to the seat such as pommels and knee-blocks. Harnesses may also be required.
Accessories
There is a wide range of accessories for wheelchairs. There are cushions, cup holders, seatbelts, storage bags, lights, and more.
A wheelchair user uses seatbelts for security or posture. Some wheelchair users want to use a seatbelt to make sure they never fall out of the wheelchair. Other wheelchair users use a seatbelt because they can not sit up straight on their own.
Notable wheelchair manufacturers
Everest and Jennings
Küschall
Otto Bock
Invacare
TiLite
Mogo Wheelchairs
| Technology | Devices | null |
21481963 | https://en.wikipedia.org/wiki/Cricket%20%28insect%29 | Cricket (insect) | Crickets are orthopteran insects which are related to bush crickets, and, more distantly, to grasshoppers. In older literature, such as Imms, "crickets" were placed at the family level (i.e. Gryllidae), but contemporary authorities including Otte now place them in the superfamily Grylloidea. The word has been used in combination to describe more distantly related taxa in the suborder Ensifera, such as king crickets and mole crickets.
Crickets have mainly cylindrically shaped bodies, round heads, and long antennae. Behind the head is a smooth, robust pronotum. The abdomen ends in a pair of long cerci; females have a long, cylindrical ovipositor. Diagnostic features include legs with 3-segmented tarsi; as with many Orthoptera, the hind legs have enlarged femora, providing power for jumping. The front wings are adapted as tough, leathery elytra, and some crickets chirp by rubbing parts of these together. The hind wings are membranous and folded when not in use for flight; many species, however, are flightless. The largest members of the family are the bull crickets, Brachytrupes, which are up to long.
Crickets are distributed all around the world except at latitudes 55° or higher, with the greatest diversity being in the tropics. They occur in varied habitats from grassland, bushes, and forests to marshes, beaches, and caves. Crickets are mainly nocturnal, and are best known for the loud, persistent, chirping song of males trying to attract females, although some species are mute. The singing species have good hearing, via the tympana on the tibiae of the front legs.
Crickets often appear as characters in literature. The Talking Cricket features in Carlo Collodi's 1883 children's book, The Adventures of Pinocchio, and in films based on the book. The insect is central to Charles Dickens's 1845 The Cricket on the Hearth and George Selden's 1960 The Cricket in Times Square. Crickets are celebrated in poems by William Wordsworth, John Keats, Du Fu and Vladimir Nazor. They are kept as pets in countries from China to Europe, sometimes for cricket fighting. Crickets are efficient at converting their food into body mass, making them a candidate for food production. They are used as human food in Southeast Asia, where they are sold deep-fried in markets as snacks. They are also used to feed carnivorous pets and zoo animals. In Brazilian folklore, crickets feature as omens of various events.
Description
Crickets are small to medium-sized insects with mostly cylindrical, somewhat vertically flattened bodies. The head is spherical with long slender antennae arising from cone-shaped scapes (first segments) and just behind these are two large compound eyes. On the forehead are three ocelli (simple eyes). The pronotum (first thoracic segment) is trapezoidal in shape, robust, and well-sclerotized. It is smooth and has neither dorsal nor lateral keels (ridges).
At the tip of the abdomen is a pair of long cerci (paired appendages on rearmost segment), and in females, the ovipositor is cylindrical, long and narrow, smooth and shiny. The femora (third segments) of the back pair of legs are greatly enlarged for jumping. The tibiae (fourth segments) of the hind legs are armed with a number of moveable spurs, the arrangement of which is characteristic of each species. The tibiae of the front legs bear one or more tympani which are used for the reception of sound.
The wings lie flat on the body and are very variable in size between species, being reduced in size in some crickets and missing in others. The fore wings are elytra made of tough chitin, acting as a protective shield for the soft parts of the body and in males, bear the stridulatory organs for the production of sound. The hind pair is membranous, folding fan-wise under the fore wings. In many species, the wings are not adapted for flight.
The largest members of the family are the -long bull crickets (Brachytrupes) which excavate burrows a metre or more deep. The tree crickets (Oecanthinae) are delicate white or pale green insects with transparent fore wings, while the field crickets (Gryllinae) are robust brown or black insects.
Distribution and habitat
Crickets have a cosmopolitan distribution, being found in all parts of the world with the exception of cold regions at latitudes higher than about 55° North and South. They have colonised many large and small islands, sometimes flying over the sea to reach these locations, or perhaps conveyed on floating timber or by human activity. The greatest diversity occurs in tropical locations, such as in Malaysia, where 88 species were heard chirping from a single location near Kuala Lumpur. A greater number than this could have been present because some species are mute.
Crickets are found in many habitats. Members of several subfamilies are found in the upper tree canopy, in bushes, and among grasses and herbs. They also occur on the ground and in caves, and some are subterranean, excavating shallow or deep burrows. Some make home in rotting wood, and certain beach-dwelling species can run and jump over the surface of water.
Biology
Defence
Crickets are relatively defenceless, soft-bodied insects. Most species are nocturnal and spend the day hidden in cracks, under bark, inside curling leaves, under stones or fallen logs, in leaf litter, or in the cracks in the ground that develop in dry weather. Some excavate their own shallow holes in rotting wood or underground and fold in their antennae to conceal their presence. Some of these burrows are temporary shelters, used for a single day, but others serve as more permanent residences and places for mating and laying eggs. Crickets burrow by loosening the soil with the mandibles and then carrying it with the limbs, flicking it backwards with the hind legs or pushing it with the head.
Other defensive strategies are the use of camouflage, fleeing, and aggression. Some species have adopted colourings, shapes, and patterns that make it difficult for predators that hunt by sight to detect them. They tend to be dull shades of brown, grey, and green that blend into their background, and desert species tend to be pale. Some species can fly, but the mode of flight tends to be clumsy, so the most usual response to danger is to scuttle away to find a hiding place. While some crickets have a weak bite, a member of the Gryllacrididae or raspy crickets from Australia were found to have the strongest bite of any insect.
Chirping
Most male crickets make a loud chirping sound by stridulation (scraping two specially textured body parts together). The stridulatory organ is located on the tegmen, or fore wing, which is leathery in texture. A large vein runs along the centre of each tegmen, with comb-like serrations on its edge forming a file-like structure, and at the rear edge of the tegmen is a scraper. The tegmina are held at an angle to the body and rhythmically raised and lowered which causes the scraper on one wing to rasp on the file on the other. The central part of the tegmen contains the "harp", an area of thick, sclerotized membrane which resonates and amplifies the volume of sound, as does the pocket of air between the tegmina and the body wall. Most female crickets lack the necessary adaptations to stridulate, so make no sound.
Several types of cricket songs are in the repertoire of some species. The calling song attracts females and repels other males, and is fairly loud. The courting song is used when a female cricket is near and encourages her to mate with the caller. A triumphal song is produced for a brief period after a successful mating and may reinforce the mating bond to encourage the female to lay some eggs rather than find another male. An aggressive song is triggered by contact chemoreceptors on the antennae that detect the presence of another male cricket.
Crickets chirp at different rates depending on their species and the temperature of their environment. Most species chirp at higher rates the higher the temperature is (about 62 chirps a minute at in one common species; each species has its own rate). The relationship between temperature and the rate of chirping is known as Dolbear's law. According to this law, counting the number of chirps produced in 14 seconds by the snowy tree cricket, common in the United States, and adding 40 will approximate the temperature in degrees Fahrenheit.
In 1975, Dr. William H. Cade discovered that the parasitic tachinid fly Ormia ochracea is attracted to the song of the cricket, and uses it to locate the male to deposit her larvae on him. It was the first known example of a natural enemy that locates its host or prey using the mating signal. Since then, many species of crickets have been found to be carrying the same parasitic fly, or related species. In response to this selective pressure, a mutation leaving males unable to chirp was observed amongst a population of Teleogryllus oceanicus on the Hawaiian island of Kauai, enabling these crickets to elude their parasitoid predators. A different mutation with the same effect was also discovered on the neighboring island of Oahu (ca. away). Recently, new "purring" males of the same species in Hawaii are able to produce a novel auditory sexual signal that can be used to attract females while greatly reducing the likelihood of parasitoid attack from the fly.
Flight
Some species, such as the ground crickets (Nemobiinae), are wingless; others have small fore wings and no hind wings (Copholandrevus), others lack hind wings and have shortened fore wings in females only, while others are macropterous, with the hind wings longer than the fore wings. In Teleogryllus, the proportion of macropterous individuals varies from very low to 100%. Probably, most species with hind wings longer than fore wings engage in flight.
Some species, such as Gryllus assimilis, take off, fly, and land efficiently and well, while other species are clumsy fliers. In some species, the hind wings are shed, leaving wing stumps, usually after dispersal of the insect by flight. In other species, they may be pulled off and consumed by the cricket itself or by another individual, probably providing a nutritional boost.
Gryllus firmus exhibits wing polymorphism; some individuals have fully functional, long hind wings and others have short wings and cannot fly. The short-winged females have smaller flight muscles, greater ovarian development, and produce more eggs, so the polymorphism adapts the cricket for either dispersal or reproduction. In some long-winged individuals, the flight muscles deteriorate during adulthood and the insect's reproductive capabilities improve.
Diet
Captive crickets are omnivorous; when deprived of their natural diet, they accept a wide range of organic foodstuffs. Some species are completely herbivorous, feeding on flowers, fruit, and leaves, with ground-based species consuming seedlings, grasses, pieces of leaf, and the shoots of young plants. Others are more predatory and include in their diet invertebrate eggs, larvae, pupae, moulting insects, scale insects, and aphids. Many are scavengers and consume various organic remains, decaying plants, seedlings, and fungi. In captivity, many species have been successfully raised on a diet of ground, commercial dry dog food, supplemented with lettuce and aphids.
Crickets have relatively powerful jaws, and several species have been known to bite humans.
Reproduction and lifecycle
Male crickets establish their dominance over each other by aggression. They start by lashing each other with their antennae and flaring their mandibles. Unless one retreats at this stage, they resort to grappling, at the same time each emitting calls that are quite unlike those uttered in other circumstances. When one achieves dominance, it sings loudly, while the loser remains silent.
Females are generally attracted to males by their calls, though in nonstridulatory species, some other mechanism must be involved. After the pair has made antennal contact, a courtship period may occur during which the character of the call changes. The female mounts the male and a single spermatophore is transferred to the external genitalia of the female. Sperm flows from this into the female's oviduct over a period of a few minutes or up to an hour, depending on species. After copulation, the female may remove or eat the spermatophore; males may attempt to prevent this with various ritualised behaviours. The female may mate on several occasions with different males.
Most crickets lay their eggs in the soil or inside the stems of plants, and to do this, female crickets have a long, needle-like or sabre-like egg-laying organ called an ovipositor. Some ground-dwelling species have dispensed with this, either depositing their eggs in an underground chamber or pushing them into the wall of a burrow. The short-tailed cricket (Anurogryllus) excavates a burrow with chambers and a defecating area, lays its eggs in a pile on a chamber floor, and after the eggs have hatched, feeds the juveniles for about a month.
Crickets are hemimetabolic insects, whose lifecycle consists of an egg stage, a larval or nymph stage that increasingly resembles the adult form as the nymph grows, and an adult stage. The egg hatches into a nymph about the size of a fruit fly. This passes through about 10 larval stages, and with each successive moult, it becomes more like an adult. After the final moult, the genitalia and wings are fully developed, but a period of maturation is needed before the cricket is ready to breed.
Inbreeding avoidance
Some species of cricket are polyandrous. In Gryllus bimaculatus, the females select and mate with multiple viable sperm donors, preferring novel mates. Female Teleogryllus oceanicus crickets from natural populations similarly mate and store sperm from multiple males. Female crickets exert a postcopulatory fertilization bias in favour of unrelated males to avoid the genetic consequences of inbreeding. Fertilization bias depends on the control of sperm transport to the sperm storage organs. The inhibition of sperm storage by female crickets can act as a form of cryptic female choice to avoid the severe negative effects of inbreeding. Controlled-breeding experiments with the cricket Gryllus firmus demonstrated inbreeding depression, as nymphal weight and early fecundity declined substantially over the generations; this was caused as expected by an increased frequency of homozygous combinations of deleterious recessive alleles.
Predators, parasites, and pathogens
Crickets have many natural enemies and are subject to various pathogens and parasites. They are eaten by large numbers of vertebrate and invertebrate predators and their hard parts are often found during the examination of animal intestines. Mediterranean house geckos (Hemidactylus turcicus) have learned that although a calling decorated cricket (Gryllodes supplicans) may be safely positioned in an out-of-reach burrow, female crickets attracted to the call can be intercepted and eaten.
The entomopathogenic fungus Metarhizium anisopliae attacks and kills crickets and has been used as the basis of control in pest populations. The insects are also affected by the cricket paralysis virus, which has caused high levels of fatalities in cricket-rearing facilities. Other fatal diseases that have been identified in mass-rearing establishments include Rickettsia and three further viruses. The diseases may spread more rapidly if the crickets become cannibalistic and eat the corpses.
Red parasitic mites sometimes attach themselves to the dorsal region of crickets and may greatly affect them. The horsehair worm Paragordius varius is an internal parasite and can control the behaviour of its cricket host and cause it to enter water, where the parasite continues its lifecycle and the cricket likely drowns. The larvae of the sarcophagid fly Sarcophaga kellyi develop inside the body cavity of field crickets. Female parasitic wasps of Rhopalosoma lay their eggs on crickets, and their developing larvae gradually devour their hosts. Other wasps in the family Scelionidae are egg parasitoids, seeking out batches of eggs laid by crickets in plant tissues in which to insert their eggs.
The fly Ormia ochracea has very acute hearing and targets calling male crickets. It locates its prey by ear and then lays its eggs nearby. The developing larvae burrow inside any crickets with which they come in contact and in the course of a week or so, devour what remains of the host before pupating. In Florida, the parasitic flies were only present in the autumn, and at that time of year, the males sang less but for longer periods. A trade-off exists for the male between attracting females and being parasitized.
Phylogeny and taxonomy
The phylogenetic relationships of the Gryllidae, summarized by Darryl Gwynne in 1995 from his own work (using mainly anatomical characteristics) and that of earlier authors, are shown in the following cladogram, with the Orthoptera divided into two main groups, Ensifera (crickets sensu lato) and Caelifera (grasshoppers). Fossil Ensifera are found from the late Carboniferous period (300 Mya) onwards, and the true crickets, Gryllidae, from the Triassic period (250 to 200 Mya).
Cladogram after Gwynne, 1995:
A phylogenetic study by Jost & Shaw in 2006 using sequences from 18S, 28S, and 16S rRNA supported the monophyly of Ensifera. Most ensiferan families were also found to be monophyletic, and the superfamily Gryllacridoidea was found to include Stenopelmatidae, Anostostomatidae, Gryllacrididae and Lezina. Schizodactylidae and Grylloidea were shown to be sister taxa, and Rhaphidophoridae and Tettigoniidae were found to be more closely related to Grylloidea than had previously been thought. The authors stated that "a high degree of conflict exists between the molecular and morphological data, possibly indicating that much homoplasy is present in Ensifera, particularly in acoustic structures." They considered that tegmen stridulation and tibial tympanae are ancestral to Ensifera and have been lost on multiple occasions, especially within the Gryllidae.
"Cricket" families
Several families and other taxa in the Ensifera may be called "crickets", including:
Within the Grylloidea
Gryllidae – "true crickets";
Mogoplistidae – scaly crickets;
Oecanthidae - tree crickets, anomalous crickets
Phalangopsidae - "spider crickets" and allies
†Protogryllidae - an extinct family
Pteroplistidae - "feather-winged crickets" of tropical Asia
Trigonidiidae - sword-tail crickets and wood or ground-crickets;
other families in the infraorder Gryllidea previously have been included:
Gryllotalpidae – mole crickets;
Myrmecophilidae – ant crickets.
Strictly, taxa in Infraorder Tettigoniidea and other superfamilies are excluded
Tettigoniidae – the bush crickets or katydids – which are quite distinct and unrelated, with 4-segmented tarsi (at least in the middle and hind legs) and females with flattened ovipositors. Also note:
within this family is the genus Anabrus – the "mormon crickets";
"bush crickets" (American usage) include members of the subfamily Trigonidiinae – which are "true crickets".
Superfamily Stenopelmatoidea – includes: king crickets (wētā), leaf-rolling, Jerusalem or sand crickets;
Superfamily Rhaphidophoroidea – cave or camel crickets;
Superfamily Schizodactyloidea - dune or splay-footed crickets.
In human culture
Folklore and myth
The folklore and mythology surrounding crickets is extensive. The singing of crickets in the folklore of Brazil and elsewhere is sometimes taken to be a sign of impending rain, or of a financial windfall. In Álvar Núñez Cabeza de Vaca's chronicles of the Spanish conquest of the Americas, the sudden chirping of a cricket heralded the sighting of land for his crew, just as their water supply had run out. In Caraguatatuba, Brazil, a black cricket in a room is said to portend illness; a grey one, money; and a green one, hope. In Alagoas state, northeast Brazil, a cricket announces death, thus it is killed if it chirps in a house. In Barbados, a loud cricket means money is coming in; hence, a cricket must not be killed or evicted if it chirps inside a house. However, another type of cricket that is less noisy forebodes illness or death.
In literature
Crickets feature as major characters in novels and children's books. Charles Dickens's 1845 novella The Cricket on the Hearth, divided into sections called "Chirps", tells the story of a cricket which chirps on the hearth and acts as a guardian angel to a family. Carlo Collodi's 1883 children's book "Le avventure di Pinocchio" (The Adventures of Pinocchio) featured "Il Grillo Parlante" (The Talking Cricket) as one of its characters. George Selden's 1960 children's book The Cricket in Times Square tells the story of Chester the cricket from Connecticut who joins a family and their other animals, and is taken to see Times Square in New York. The story, which won the Newbery Honor, came to Selden on hearing a real cricket chirp in Times Square.
Souvenirs entomologiques, a book written by the French entomologist Jean-Henri Fabre, devotes a whole chapter to the cricket, discussing its construction of a burrow and its song-making. The account is mainly of the field cricket, but also mentions the Italian cricket.
Crickets have from time to time appeared in poetry. William Wordsworth's 1805 poem The Cottager to Her Infant includes the couplet "The kitten sleeps upon the hearth, The crickets long have ceased their mirth". John Keats's 1819 poem Ode to Autumn includes the lines "Hedge-crickets sing; and now with treble soft / The redbreast whistles from a garden-croft". The Chinese Tang dynasty poet Du Fu (712–770) wrote a poem that in the translation by J. P. Seaton begins "House cricket ... Trifling thing. And yet how his mournful song moves us. Out in the grass his cry was a tremble, But now, he trills beneath our bed, to share his sorrow."
As pets and fighting animals
Crickets are kept as pets and are considered good luck in some countries; in China, they are sometimes kept in cages or in hollowed-out gourds specially created in novel shapes. The practice was common in Japan for thousands of years; it peaked in the 19th century, though crickets are still sold at pet shops. It is also common to have them as caged pets in some European countries, particularly in the Iberian Peninsula. Cricket fighting is a traditional Chinese pastime that dates back to the Tang dynasty (618–907). Originally an indulgence of emperors, cricket fighting later became popular among commoners. The dominance and fighting ability of males does not depend on strength alone; it has been found that they become more aggressive after certain pre-fight experiences such as isolation, or when defending a refuge. Crickets forced to fly for a short while will afterwards fight for two to three times longer than they otherwise would.
As food
In the southern part of Asia, including Cambodia, Laos, and Thailand, crickets commonly are eaten as a snack, prepared by deep frying soaked and cleaned insects. In Thailand, there are 20,000 farmers rearing crickets, with an estimated production of 7,500 tons per year. The United Nations' Food and Agriculture Organization has implemented a project in Laos to improve cricket farming and, consequently, food security. The food conversion efficiency of house crickets (Acheta domesticus) is 1.7, some five times higher than that for beef cattle, and if their fecundity is taken into account, 15 to 20 times higher.
Cricket flour may be used as an additive to consumer foods such as pasta, bread, crackers, and cookies. Cricket flour is used in protein bars, pet foods, livestock feed, nutraceuticals, and other industrial applications. The United Nations says that the use of insect protein, such as cricket flour, could be critical in feeding the growing population of the planet while being less damaging to the environment.
Crickets are also raised as food for carnivorous zoo animals, laboratory animals, and pets. They may be "gut loaded" with additional minerals, such as calcium, to provide a balanced diet for predators such as tree frogs (Hylidae).
Common expressions
By the 19th century "cricket" and "crickets" were in use as euphemisms for using Christ as an interjection. The addition of "Jiminy" (a variation of "Gemini"), sometimes shortened to "Jimmy" created the expressions "Jiminy Cricket!" or "Jimmy Crickets!" as less blasphemous alternatives to exclaiming "Jesus Christ!"
By the end of the 20th century the sound of chirping crickets came to represent quietude in literature, theatre and film. From this sentiment arose expressions equating "crickets" with silence altogether, particularly when a group of assembled people makes no noise. These expressions have grown from the more descriptive, "so quiet that you can hear crickets," to simply saying, "crickets" as shorthand for "complete silence."
In popular culture
Cricket characters feature in the Walt Disney animated movies Pinocchio (1940), where Jiminy Cricket becomes the title character's conscience, and in Mulan (1998), where Cri-Kee is carried in a cage as a symbol of luck, in the Asian manner. The Crickets was the name of Buddy Holly's rock and roll band; Holly's home town baseball team in the 1990s was called the Lubbock Crickets. Cricket is the name of a US children's literary magazine founded in 1973; it uses a cast of insect characters.
| Biology and health sciences | Orthoptera | null |
3857434 | https://en.wikipedia.org/wiki/Scansoriopterygidae | Scansoriopterygidae | Scansoriopterygidae (meaning "climbing wings") is an extinct family of climbing and gliding maniraptoran dinosaurs. Scansoriopterygids are known from five well-preserved fossils, representing four species, unearthed in the Tiaojishan Formation fossil beds (dating to the mid-late Jurassic Period) of Liaoning and Hebei, China.
Scansoriopteryx heilmanni (and its likely synonym Epidendrosaurus ninchengensis) was the first non-avian dinosaur found that had clear adaptations to an arboreal or semi-arboreal lifestyle–it is likely that they spent much of their time in trees. Both specimens showed features indicating they were juveniles, which made it difficult to determine their exact relationship to other non-avian dinosaurs and birds. It was not until the description of Epidexipteryx hui in 2008 that a subadult specimen was known. In 2015, the discovery of an adult specimen belonging to the species Yi qi showed that scansoriopterygids were not only climbers but also had adaptations that could have been used for gliding flight. The recently discovered (in 2019) Ambopteryx also supports this. The earlier described Pedopenna may also be a member of this clade.
Description
Scansoriopterygids are among the smallest non-avian dinosaurs known. The juvenile specimens of Scansoriopteryx are the size of house sparrows, about long, while the subadult type specimen of Epidexipteryx is about the size of a pigeon, about long (not including the tail feathers).
Scansoriopterygids differentiate from other theropod dinosaurs in part by their extremely long third fingers, which were longer than the first and second digits of the hand. In all other known theropods, the second finger is the longest. At least two species, Yi and Ambopteryx, also had a long "styliform" bone growing from the wrist, which, along with the third finger, helped support a bat-like wing membrane used for gliding. This use of a long finger to support a wing membrane is only superficially similar to the wing arrangement in pterosaurs, even though it is physically more bat-like.
Other features shared within the group include short and high skulls with down turned lower jaws and large front teeth, and long arms. Tail length, however, varied significantly among scansoriopterygids. Epidexipteryx had a short tail (70% the length of the torso), anchoring long tail feathers, while Scansoriopteryx had a very long tail (over three times as long as the torso) with a short spray of feathers at the tip. All three described scansoripterygid specimens preserve the fossilized traces of feathers covering their bodies.
Classification
Scansoriopterygidae was created as a family-level taxon by Stephen Czerkas and Yuan Chongxi in 2002. Some scientists, such as Paul Sereno, initially considered the concept redundant because the group was originally monotypic, containing only the single genus and species Scansoriopteryx heilmanni. Additionally, the group lacked a phylogenetic definition. However, in 2008 Zhang et al. reported another scansoriopterygid, Epidexipteryx, and defined Scansoriopterygidae as a clade comprising most recent common ancestor of Epidexipteryx and Epidendrosaurus (=Scansoriopteryx) plus all its descendants.
The exact taxonomic placement of this group was initially uncertain and controversial. When describing the first validly published specimen in 2002 (Scansoriopteryx heilmanni), Czerkas and Yuan proposed that various primitive features of the skeleton (including a primitive, "saurischian-style" pubis and primitive hip joint) showed that scansoriopterygids, along with other maniraptorans and birds, split from other theropods very early in dinosaur evolution. However, this interpretation has not been followed by most other researchers. In a 2007 cladistic analysis of relationships among coelurosaurs, Phil Senter found Scansoriopteryx to be a member of the clade Avialae. This view was supported by a second phylogenetic analysis performed by Zhang et al. in 2008.
A subsequent phylogenetic analysis conducted by Agnolín and Novas (2011) recovered scansoriopterygids not as avialans, but as basal members of the clade Paraves remaining in unresolved polytomy with alvarezsaurids and the clade Eumaniraptora (containing avialans and deinonychosaurs).
Turner, Makovicky and Norell (2012) included only Epidexipteryx hui in their primary phylogenetic analysis, as a full-grown specimen of this species is known; regarding Scansoriopteryx/Epidendrosaurus, the authors were worried that including it in the primary analysis would be problematic, because it is only known from juvenile specimens, which "do not necessarily preserve all the adult morphology needed to accurately place a taxon phylogenetically" (Turner, Makovicky and Norell 2012, p. 89). Epidexipteryx was recovered as basal paravian that didn't belong to Eumaniraptora. The authors did note that its phylogenetic position is unstable; constraining Epidexipteryx hui as a basal avialan required two additional steps compared to the most parsimonious solution, while constraining it as a basal member of Oviraptorosauria required only one additional step. A separate exploratory analysis included Scansoriopteryx/Epidendrosaurus, which was recovered as a basal member of Avialae; the authors noted that it did not clade with Epidexipteryx, which stayed outside Eumaniraptora. Constraining the monophyly of Scansoriopterygidae required four additional steps and moved Epidexipteryx into Avialae.
A monophyletic Scansoriopterygidae was recovered by Godefroit et al. (2013); the authors found scansoriopterygids to be basalmost members of Paraves and the sister group to the clade containing Avialae and Deinonychosauria. Agnolín and Novas (2013) recovered scansoriopterygids as non-paravian maniraptorans and the sister group to Oviraptorosauria. Brusatte et al. (2014) also found Epidexipteryx to be a basal oviraptorosaur along with Pedopenna. The Bayesian analysis of Cau (2018) placed scansoriopterygids in Oviraptorosauria again, while the parsimony analysis placed them in the base of Avialae, and included Xiaotingia in Scansoriopterygidae as sister to the rest of the group. Pittman et al. (2020) again found scansoriopterygids to be basal oviraptorosaurs. Cau (2024) placed Scansoriopterygidae within Anchiornithidae.
The cladogram below follows the results of a phylogenetic study by Lefèvre et al., 2014:
Sorkin (2021) argues scansoriopterygids provide evidence that all pennaraptorans evolved from scansorial gliding ancestors.
Provenance and paleoecology
The fossil remains of Epidexipteryx, Scansoriopteryx and Yi were all recovered from the Tiaojishan Formation of northeastern China, and the former two were specifically found in the Daohugou Beds. A study published in 2008 refined the possible age range of this formation, finding that the lower boundary of the Tiaojishan was formed 165 Ma ago, and the upper boundary somewhere between 156 and 153 Ma ago.
The known scansoriopterygids of the Daohugou biota inhabited a humid, temperate forest made up of a variety of prehistoric trees including species of ginkgo and conifer. The understory would have been dominated by plants such as club mosses, horsetails, cycads, and ferns.
The scansoriopterygids would have lived alongside synapsids such as the aquatic Castorocauda, arboreal gliding mammal Volaticotherium and various types of gliding haramiyidans, the rhamphorhynchoid pterosaurs Jeholopterus and Pterorhynchus, as well as a diverse range of insect life (including mayflies and beetles) and several species of salamander.
Paleobiology
Climbing
In the initial descriptions of the first two scansoriopterygid specimens, scientists studying these animals used several lines of evidence to argue that they were arboreal (tree-climbing), and the first known non-avian dinosaurs with clear climbing adaptations.
Zhang and colleagues considered Scansoriopteryx to be arboreal based on the elongated nature of the hand and specializations of the foot. These authors stated that the long hand and strongly curved claws were adaptations for climbing and moving around among tree branches. They viewed this as an early stage in the evolution of the bird wing, stating that the forelimbs became well-developed for climbing, and that this development later led to the evolution of a wing capable of flight. They argued that long, grasping hands are more suited to climbing than to flight, since most flying birds have relatively short hands. Zhang et al. also noted that the foot of Scansoriopteryx is unique among non-avian theropods; while Scansoriopteryx does not preserve a reversed hallux (the backward-facing toe seen in modern perching birds), its foot was very similar in construction to primitive perching birds like Cathayornis and Longipteryx. These adaptations for grasping ability in all four limbs, the authors argued, makes it likely that Scansoriopteryx spent a significant amount of time living in trees.
In describing Scansoriopteryx, Czerkas and Yuan also described evidence for an arboreal lifestyle. They noted that, unlike all modern bird hatchlings, the forelimbs of Scansoriopteryx are longer than the hind limbs. The authors argued that this anomaly indicates the forelimbs played an important role in locomotion even at an extremely early developmental stage. Scansoriopteryx has a better-preserved foot than the type of Epidendrosaurus, and the authors interpreted the hallux as reversed, the condition of a backward-pointing toe being widespread among modern tree-dwelling birds. Furthermore, the authors pointed to the stiffened tail of Scansoriopteryx as a tree-climbing adaptation. The tail may have been used as a prop, much like the tails of modern woodpeckers. Comparison with the hands of modern climbing species with elongated third digits, like iguanid lizards, also supports the tree-climbing hypothesis. Indeed, the hands of Scansoriopteryx are much better adapted to climbing than the modern tree-climbing hatchling of the hoatzin.
Feathers
Both juvenile scansoriopterygid specimens preserve impressions of simple, down-like feathers, especially around the hand and arm. The longer feathers in this region led Czerkas and Yuan to speculate that adult scansoriopterygids may have had reasonably well-developed wing feathers which could have aided in leaping or rudimentary gliding, though they ruled out the possibility that Scansoriopteryx could have achieved powered flight. Like other maniraptorans, scansoriopterygids had a semilunate carpal (half-moon shaped wrist bone) that allowed for bird-like folding motion in the hand. Even if powered flight was not possible, this motion could have aided maneuverability in leaping from branch to branch.
The adult specimen of Epidexipteryx lacked preserved feathers around the forelimbs, but preserved simple feathers on the body and long, ribbon-like feathers on the tail. The tail feathers, likely used in display, consisted of a central shaft (rachis) and unbranched vane (unlike the vanes of modern feathers, which are broken up into smaller filaments or barbs).
Yi also preserves feathers. These are notably very simple for a member of Pennaraptora (a clade of which scansoriopterygids are usually considered members), being "paintbrush-like", with long quill-like bases topped by sprays of thinner filaments. The feathers covered most of the body, starting near the tip of the snout. The head and neck feathers were long and formed a thick coat, and the body feathers were even longer and denser, making it difficult for scientists to study their detailed structure.
Gliding membranes
At least two species, Yi qi and Ambopteryx longibrachium, developed a patagium, supporting it with the elongated third finger as well as a unique styliform wrist bone akin to similar structures in flying squirrels, bats, pterosaurs and anomalures. Though propatagia are known in birds and similar dinosaurs, scansoriopterygids were the only known dinosaurs to develop true membranous wings, most notably so due to the presence of already fairly derived feathers.
Prior to the discovery of Yi, Italian palaeontologist Andrea Cau had informally suggested that membranes may have been present in Scansoriopteryx, supported by its elongated third finger, due to their similarity to the wing fingers of pterosaurs, a hypothesis he later also applied to Epidexipteryx.
| Biology and health sciences | Theropods | Animals |
3861311 | https://en.wikipedia.org/wiki/Pupfish | Pupfish | Pupfish are a group of small killifish belonging to ten genera of the family Cyprinodontidae of ray-finned fish. Pupfish are especially noted for being found in extreme and isolated situations. They are primarily found in North America, South America, and the Caribbean region. As of August 2006, 120 nominal species and 9 subspecies were known. Several pupfish species are extinct and most extant species are listed. In the U.S., the most well-known pupfish species may be the Devils Hole pupfish, native to Devils Hole on the Nevada side of Death Valley National Park. Since 1995 the Devils Hole pupfish has been in a nearly steady decline, where it was close to extinction at 35–68 fish in 2013.
The common name is said to derive from the mating habits of the males, whose activities vaguely resemble puppies at play; Carl L. Hubbs, a prominent ichthyologist and one of the first people to take an interest in them, coined the name after he observed their "playful" circling and tussling, which is actually the aggressive behavior of territorial males.
In spite of their name, the cyprinodontids are not closely related to Cyprinidae, or carp family. They were formerly considered near allies of the pikes and their relatives, as they share some features: a flat head with protractile mouth beset with cardiform, villiform, or compressed, bi- or tricuspid teeth, generally large scales, and the absence of a well-developed lateral line. However, they are now generally assigned to the order Cyprinodontiformes. Several forms occur in the fossil records of the Oligocene and Miocene beds of Europe. Pupfish from San Salvador Island were able to diversify into multiple species with different eating habits due to interbreeding with pupfish from other islands, mainly Caribbean.
Most pupfish are inhabitants of fresh and brackish waters. Many species are ovoviviparous; often the sexes are dissimilar, the female being larger and less brilliantly coloured, with smaller fins; the anal fin of the male may be modified into an intromittent organ by means of which internal fertilization takes place. Most pupfishes' diet consists, mainly, of algae, decaying vegetation, and any insects they can get.
Genera
Cualac Miller, 1956
Cubanichthys C.L. Hubbs, 1926
Cyprinodon Lacepède, 1803
Floridichthys C.L. Hubbs, 1926
Garmanella C.L. Hubbs, 1936
Jordanella Goode & Bean, 1879
†Megupsilon Miller & Walters, 1972 (extinct 2014)
Orestias Valenciennes 1839
Pseudorestias Arratia, Vila, Lam, Guerrero & Quezada-Romegialli, 2017
Yssolebias Huber 2012 (possibly extinct, only known from one old specimen)
Pupfish evolution
Pupfish on the island of San Salvador, Bahamas, have a large adaptive diversification in only two small lakes. They evolve 50–130 times faster than any other species of pupfish. This is also the fastest morphological diversification seen in any fish that has been documented. It is believed that this diversification is because of their ecological niches.
Three species of pupfish on the island of San Salvador, Bahamas, all live in salty lakes. These pupfish are able to take advantage of different food sources so they can all coexist. One species feeds on only the scales of other pupfish. Another has a modified jaw to be able to eat snails and ostracods.
Before the 1990s, Lake Chichancanab, in Mexico, was full of brackish water and another five species of pupfish were found there. Cyprinodon maya was the largest pupfish, and it ate other fish. Cyprinodon simus was the second smallest, and it ate zooplankton. These species are now considered extinct in the wild because of an invasive species of African tilapia.
The Death Valley pupfish evolve 5–10 faster than average and are known for their abilities to survive in extremely hot waters. Cyprinodon diabolis eat algae off a rock shelf near the surface of the deep pool they live in.
North American pupfish
The pupfish found in Death Valley were once thought to be one main species. They were once all found in Lake Manly, a glacial lake over 620 square miles (1,600 km2), roughly 185,000-128,000 years ago. Over time this lake dried up and started to separate into smaller lakes or ponds. As this drying happened the pupfish became separated into different ponds and started to divergently evolve. There are thought to be two main subspecies of Death Valley pupfish (C. salinus and C. milleri) present. These are both considered endangered since they are only found in one area of the world. Cyprinodon pachycephalus live in extremely hot waters, 114 °F (45.5 °C).
The Devils Hole pupfish (Cyprinodon diabolis) is a specific species native to Nevada. There are fewer than 200 individuals since 2005. Their population size usually fluctuates between 37 and 400 fish. They are considered one of the world's rarest fish. These fish live in 94 °F (34.4 °C) waters.
| Biology and health sciences | Acanthomorpha | Animals |
34578727 | https://en.wikipedia.org/wiki/Variable%20renewable%20energy | Variable renewable energy | Variable renewable energy (VRE) or intermittent renewable energy sources (IRES) are renewable energy sources that are not dispatchable due to their fluctuating nature, such as wind power and solar power, as opposed to controllable renewable energy sources, such as dammed hydroelectricity or bioenergy, or relatively constant sources, such as geothermal power.
The use of small amounts of intermittent power has little effect on grid operations. Using larger amounts of intermittent power may require upgrades or even a redesign of the grid infrastructure.
Options to absorb large shares of variable energy into the grid include using storage, improved interconnection between different variable sources to smooth out supply, using dispatchable energy sources such as hydroelectricity and having overcapacity, so that sufficient energy is produced even when weather is less favourable. More connections between the energy sector and the building, transport and industrial sectors may also help.
Background and terminology
The penetration of intermittent renewables in most power grids is low: global electricity generation in 2021 was 7% wind and 4% solar. However, in 2021 Denmark, Luxembourg and Uruguay generated over 40% of their electricity from wind and solar. Characteristics of variable renewables include their unpredictability, variability, and low operating costs. These, along with renewables typically being asynchronous generators, provide a challenge to grid operators, who must make sure supply and demand are matched. Solutions include energy storage, demand response, availability of overcapacity and sector coupling. Smaller isolated grids may be less tolerant to high levels of penetration.
Matching power demand to supply is not a problem specific to intermittent power sources. Existing power grids already contain elements of uncertainty including sudden and large changes in demand and unforeseen power plant failures. Though power grids are already designed to have some capacity in excess of projected peak demand to deal with these problems, significant upgrades may be required to accommodate large amounts of intermittent power.
Several key terms are useful for understanding the issue of intermittent power sources. These terms are not standardized, and variations may be used. Most of these terms also apply to traditional power plants.
Intermittency or variability is the extent to which a power source fluctuates. This has two aspects: a predictable variability, such as the day-night cycle, and an unpredictable part (imperfect local weather forecasting). The term intermittent can be used to refer to the unpredictable part, with variable then referring to the predictable part.
Dispatchability is the ability of a given power source to add output on demand. The concept is distinct from intermittency; dispatchability is one of several ways system operators match supply (generator's output) to system demand (technical loads).
Penetration is the amount of electricity generated from a particular source as a percentage of annual consumption.
Nominal power or nameplate capacity is the theoretical output registered with authorities for classifying the unit. For intermittent power sources, such as wind and solar, nameplate power is the source's output under ideal conditions, such as maximum usable wind or high sun on a clear summer day.
Capacity factor, average capacity factor, or load factor is the ratio of actual electrical generation over a given period of time, usually a year, to actual generation in that time period. Basically, it is the ratio between the how much electricity a plant produced and how much electricity a plant would have produced if were running at its nameplate capacity for the entire time period.
Firm capacity or firm power is "guaranteed by the supplier to be available at all times during a period covered by a commitment".
Capacity credit: the amount of conventional (dispatchable) generation power that can be potentially removed from the system while keeping the reliability, usually expressed as a percentage of the nominal power.
Foreseeability or predictability is how accurately the operator can anticipate the generation: for example tidal power varies with the tides but is completely foreseeable because the orbit of the moon can be predicted exactly, and improved weather forecasts can make wind power more predictable.
Sources
Dammed hydroelectricity, biomass and geothermal are dispatchable as each has a store of potential energy; wind and solar without storage can be decreased (curtailed) but are not dispatchable.
Wind power
Grid operators use day ahead forecasting to determine which of the available power sources to use the next day, and weather forecasting is used to predict the likely wind power and solar power output available. Although wind power forecasts have been used operationally for decades, the IEA is organizing international collaboration to further improve their accuracy.
Wind-generated power is a variable resource, and the amount of electricity produced at any given point in time by a given plant will depend on wind speeds, air density, and turbine characteristics, among other factors. If wind speed is too low then the wind turbines will not be able to make electricity, and if it is too high the turbines will have to be shut down to avoid damage. While the output from a single turbine can vary greatly and rapidly as local wind speeds vary, as more turbines are connected over larger and larger areas the average power output becomes less variable.
Intermittence: Regions smaller than synoptic scale, less than about 1000 km long, the size of an average country, have mostly the same weather and thus around the same wind power, unless local conditions favor special winds. Some studies show that wind farms spread over a geographically diverse area will as a whole rarely stop producing power altogether. This is rarely the case for smaller areas with uniform geography such as Ireland, Scotland and Denmark which have several days per year with little wind power.
Capacity factor: Wind power typically has an annual capacity factor of 25–50%, with offshore wind outperforming onshore wind.
Dispatchability: Because wind power is not by itself dispatchable wind farms are sometimes built with storage.
Capacity credit: At low levels of penetration, the capacity credit of wind is about the same as the capacity factor. As the concentration of wind power on the grid rises, the capacity credit percentage drops.
Variability: Site dependent. Sea breezes are much more constant than land breezes. Seasonal variability may reduce output by 50%.
Reliability: A wind farm has high technical reliability when the wind blows. That is, the output at any given time will only vary gradually due to falling wind speeds or storms, the latter necessitating shut downs. A typical wind farm is unlikely to have to shut down in less than half an hour at the extreme, whereas an equivalent-sized power station can fail totally instantaneously and without warning. The total shutdown of wind turbines is predictable via weather forecasting. The average availability of a wind turbine is 98%, and when a turbine fails or is shut down for maintenance it only affects a small percentage of the output of a large wind farm.
Predictability: Although wind is variable, it is also predictable in the short term. There is an 80% chance that wind output will change less than 10% in an hour and a 40% chance that it will change 10% or more in 5 hours.
Because wind power is generated by large numbers of small generators, individual failures do not have large impacts on power grids. This feature of wind has been referred to as resiliency.
Solar power
Intermittency inherently affects solar energy, as the production of renewable electricity from solar sources depends on the amount of sunlight at a given place and time. Solar output varies throughout the day and through the seasons, and is affected by dust, fog, cloud cover, frost or snow. Many of the seasonal factors are fairly predictable, and some solar thermal systems make use of heat storage to produce grid power for a full day.
Variability: In the absence of an energy storage system, solar does not produce power at night, little in bad weather and varies between seasons. In many countries, solar produces most energy in seasons with low wind availability and vice versa.
Capacity factor Standard photovoltaic solar has an annual average capacity factor of 10-20%, but panels that move and track the sun have a capacity factor up to 30%. Thermal solar parabolic trough with storage 56%. Thermal solar power tower with storage 73%.
The impact of intermittency of solar-generated electricity will depend on the correlation of generation with demand. For example, solar thermal power plants such as Nevada Solar One are somewhat matched to summer peak loads in areas with significant cooling demands, such as the south-western United States. Thermal energy storage systems like the small Spanish Gemasolar Thermosolar Plant can improve the match between solar supply and local consumption. The improved capacity factor using thermal storage represents a decrease in maximum capacity, and extends the total time the system generates power.
Run-of-the-river hydroelectricity
In many countries new large dams are no longer being built, because of the environmental impact of reservoirs. Run of the river projects have continued to be built. The absence of a reservoir results in both seasonal and annual variations in electricity generated.
Tidal power
Tidal power is the most predictable of all the variable renewable energy sources. The tides reverse twice a day, but they are never intermittent, on the contrary they are completely reliable.
Wave power
Waves are primarily created by wind, so the power available from waves tends to follow that available from wind, but due to the mass of the water is less variable than wind power. Wind power is proportional to the cube of the wind speed, while wave power is proportional to the square of the wave height.
Solutions for their integration
The displaced dispatchable generation could be coal, natural gas, biomass, nuclear, geothermal or storage hydro. Rather than starting and stopping nuclear or geothermal, it is cheaper to use them as constant base load power. Any power generated in excess of demand can displace heating fuels, be converted to storage or sold to another grid. Biofuels and conventional hydro can be saved for later when intermittents are not generating power. Some forecast that “near-firm” renewables (batteries with solar and/or wind) power will be cheaper than existing nuclear by the late 2020s: therefore they say base load power will not be needed.
Alternatives to burning coal and natural gas which produce fewer greenhouse gases may eventually make fossil fuels a stranded asset that is left in the ground. Highly integrated grids favor flexibility and performance over cost, resulting in more plants that operate for fewer hours and lower capacity factors.
All sources of electrical power have some degree of variability, as do demand patterns which routinely drive large swings in the amount of electricity that suppliers feed into the grid. Wherever possible, grid operations procedure are designed to match supply with demand at high levels of reliability, and the tools to influence supply and demand are well-developed. The introduction of large amounts of highly variable power generation may require changes to existing procedures and additional investments.
The capacity of a reliable renewable power supply, can be fulfilled by the use of backup or extra infrastructure and technology, using mixed renewables to produce electricity above the intermittent average, which may be used to meet regular and unanticipated supply demands. Additionally, the storage of energy to fill the shortfall intermittency or for emergencies can be part of a reliable power supply.
In practice, as the power output from wind varies, partially loaded conventional plants, which are already present to provide response and reserve, adjust their output to compensate. While low penetrations of intermittent power may use existing levels of response and spinning reserve, the larger overall variations at higher penetrations levels will require additional reserves or other means of compensation.
Operational reserve
All managed grids already have existing operational and "spinning" reserve to compensate for existing uncertainties in the power grid. The addition of intermittent resources such as wind does not require 100% "back-up" because operating reserves and balancing requirements are calculated on a system-wide basis, and not dedicated to a specific generating plant.
Some gas, or hydro power plants are partially loaded and then controlled to change as demand changes or to replace rapidly lost generation. The ability to change as demand changes is termed "response". The ability to quickly replace lost generation, typically within timescales of 30 seconds to 30 minutes, is termed "spinning reserve".
Generally thermal plants running as peaking plants will be less efficient than if they were running as base load. Hydroelectric facilities with storage capacity, such as the traditional dam configuration, may be operated as base load or peaking plants.
Grids can contract for grid battery plants, which provide immediately available power for an hour or so, which gives time for other generators to be started up in the event of a failure, and greatly reduces the amount of spinning reserve required.
Demand response
Demand response is a change in consumption of energy to better align with supply. It can take the form of switching off loads, or absorb additional energy to correct supply/demand imbalances. Incentives have been widely created in the American, British and French systems for the use of these systems, such as favorable rates or capital cost assistance, encouraging consumers with large loads to take them offline whenever there is a shortage of capacity, or conversely to increase load when there is a surplus.
Certain types of load control allow the power company to turn loads off remotely if insufficient power is available. In France large users such as CERN cut power usage as required by the System Operator - EDF under the encouragement of the EJP tariff.
Energy demand management refers to incentives to adjust use of electricity, such as higher rates during peak hours. Real-time variable electricity pricing can encourage users to adjust usage to take advantage of periods when power is cheaply available and avoid periods when it is more scarce and expensive. Some loads such as desalination plants, electric boilers and industrial refrigeration units, are able to store their output (water and heat). Several papers also concluded that Bitcoin mining loads would reduce curtailment, hedge electricity price risk, stabilize the grid, increase the profitability of renewable energy power stations and therefore accelerate transition to sustainable energy. But others argue that Bitcoin mining can never be sustainable.
Instantaneous demand reduction. Most large systems also have a category of loads which instantly disconnect when there is a generation shortage, under some mutually beneficial contract. This can give instant load reductions or increases.
Storage
At times of low load where non-dispatchable output from wind and solar may be high, grid stability requires lowering the output of various dispatchable generating sources or even increasing controllable loads, possibly by using energy storage to time-shift output to times of higher demand. Such mechanisms can include:
Pumped storage hydropower is the most prevalent existing technology used, and can substantially improve the economics of wind power. The availability of hydropower sites suitable for storage will vary from grid to grid. Typical round trip efficiency is 80%.
Traditional lithium-ion is the most common type used for grid-scale battery storage . Rechargeable flow batteries can serve as a large capacity, rapid-response storage medium. Hydrogen can be created through electrolysis and stored for later use.
Flywheel energy storage systems have some advantages over chemical batteries. Along with substantial durability which allows them to be cycled frequently without noticeable life reduction, they also have very fast response and ramp rates. They can go from full discharge to full charge within a few seconds. They can be manufactured using non-toxic and environmentally friendly materials, easily recyclable once the service life is over.
Thermal energy storage stores heat. Stored heat can be used directly for heating needs or converted into electricity. In the context of a CHP plant a heat storage can serve as a functional electricity storage at comparably low costs. Ice storage air conditioning Ice can be stored inter seasonally and can be used as a source of air-conditioning during periods of high demand. Present systems only need to store ice for a few hours but are well developed.
Storage of electrical energy results in some lost energy because storage and retrieval are not perfectly efficient. Storage also requires capital investment and space for storage facilities.
Geographic diversity and complementing technologies
The variability of production from a single wind turbine can be high. Combining any additional number of turbines, for example, in a wind farm, results in lower statistical variation, as long as the correlation between the output of each turbine is imperfect, and the correlations are always imperfect due to the distance between each turbine. Similarly, geographically distant wind turbines or wind farms have lower correlations, reducing overall variability. Since wind power is dependent on weather systems, there is a limit to the benefit of this geographic diversity for any power system.
Multiple wind farms spread over a wide geographic area and gridded together produce power more constantly and with less variability than smaller installations. Wind output can be predicted with some degree of confidence using weather forecasts, especially from large numbers of turbines/farms. The ability to predict wind output is expected to increase over time as data is collected, especially from newer facilities.
Electricity produced from solar energy tends to counterbalance the fluctuating supplies generated from wind. Normally it is windiest at night and during cloudy or stormy weather, and there is more sunshine on clear days with less wind. Besides, wind energy has often a peak in the winter season, whereas solar energy has a peak in the summer season; the combination of wind and solar reduces the need for dispatchable backup power.
In some locations, electricity demand may have a high correlation with wind output, particularly in locations where cold temperatures drive electric consumption, as cold air is denser and carries more energy.
The allowable penetration may be increased with further investment in standby generation. For instance some days could produce 80% intermittent wind and on the many windless days substitute 80% dispatchable power like natural gas, biomass and Hydro.
Areas with existing high levels of hydroelectric generation may ramp up or down to incorporate substantial amounts of wind. Norway, Brazil, and Manitoba all have high levels of hydroelectric generation, Quebec produces over 90% of its electricity from hydropower, and Hydro-Québec is the largest hydropower producer in the world. The U.S. Pacific Northwest has been identified as another region where wind energy is complemented well by existing hydropower. Storage capacity in hydropower facilities will be limited by size of reservoir, and environmental and other considerations.
Connecting grid internationally
It is often feasible to export energy to neighboring grids at times of surplus, and import energy when needed. This practice is common in Europe and between the US and Canada. Integration with other grids can lower the effective concentration of variable power: for instance, Denmark's high penetration of VRE, in the context of the German/Dutch/Scandinavian grids with which it has interconnections, is considerably lower as a proportion of the total system. Hydroelectricity that compensates for variability can be used across countries.
The capacity of power transmission infrastructure may have to be substantially upgraded to support export/import plans. Some energy is lost in transmission. The economic value of exporting variable power depends in part on the ability of the exporting grid to provide the importing grid with useful power at useful times for an attractive price.
Sector coupling
Demand and generation can be better matched when sectors such as mobility, heat and gas are coupled with the power system. The electric vehicle market is for instance expected to become the largest source of storage capacity. This may be a more expensive option appropriate for high penetration of variable renewables, compared to other sources of flexibility. The International Energy Agency says that sector coupling is needed to compensate for the mismatch between seasonal demand and supply.
Electric vehicles can be charged during periods of low demand and high production, and in some places send power back from the vehicle-to-grid.
Penetration
Penetration refers to the proportion of a primary energy (PE) source in an electric power system, expressed as a percentage. There are several methods of calculation yielding different penetrations. The penetration can be calculated either as:
the nominal capacity (installed power) of a PE source divided by the peak load within an electric power system; or
the nominal capacity (installed power) of a PE source divided by the total capacity of the electric power system; or
the electrical energy generated by a PE source in a given period, divided by the demand of the electric power system in this period.
The level of penetration of intermittent variable sources is significant for the following reasons:
Power grids with significant amounts of dispatchable pumped storage, hydropower with reservoir or pondage or other peaking power plants such as natural gas-fired power plants are capable of accommodating fluctuations from intermittent power more easily.
Relatively small electric power systems without strong interconnection (such as remote islands) may retain some existing diesel generators but consuming less fuel, for flexibility until cleaner energy sources or storage such as pumped hydro or batteries become cost-effective.
In the early 2020s wind and solar produce 10% of the world's electricity, but supply in the 40-55% penetration range has already been implemented in several systems, with over 65% planned for the UK by 2030.
There is no generally accepted maximum level of penetration, as each system's capacity to compensate for intermittency differs, and the systems themselves will change over time. Discussion of acceptable or unacceptable penetration figures should be treated and used with caution, as the relevance or significance will be highly dependent on local factors, grid structure and management, and existing generation capacity.
For most systems worldwide, existing penetration levels are significantly lower than practical or theoretical maximums.
Maximum penetration limits
Maximum penetration of combined wind and solar is estimated at around 70% to 90% without regional aggregation, demand management or storage; and up to 94% with 12 hours of storage. Economic efficiency and cost considerations are more likely to dominate as critical factors; technical solutions may allow higher penetration levels to be considered in future, particularly if cost considerations are secondary.
Economic impacts of variability
Estimates of the cost of wind and solar energy may include estimates of the "external" costs of wind and solar variability, or be limited to the cost of production. All electrical plant has costs that are separate from the cost of production, including, for example, the cost of any necessary transmission capacity or reserve capacity in case of loss of generating capacity. Many types of generation, particularly fossil fuel derived, will have cost externalities such as pollution, greenhouse gas emission, and habitat destruction, which are generally not directly accounted for.
The magnitude of the economic impacts is debated and will vary by location, but is expected to rise with higher penetration levels. At low penetration levels, costs such as operating reserve and balancing costs are believed to be insignificant.
Intermittency may introduce additional costs that are distinct from or of a different magnitude than for traditional generation types. These may include:
Transmission capacity: transmission capacity may be more expensive than for nuclear and coal generating capacity due to lower load factors. Transmission capacity will generally be sized to projected peak output, but average capacity for wind will be significantly lower, raising cost per unit of energy actually transmitted. However transmission costs are a low fraction of total energy costs.
Additional operating reserve: if additional wind and solar does not correspond to demand patterns, additional operating reserve may be required compared to other generating types, however this does not result in higher capital costs for additional plants since this is merely existing plants running at low output - spinning reserve. Contrary to statements that all wind must be backed by an equal amount of "back-up capacity", intermittent generators contribute to base capacity "as long as there is some probability of output during peak periods". Back-up capacity is not attributed to individual generators, as back-up or operating reserve "only have meaning at the system level".
Balancing costs: to maintain grid stability, some additional costs may be incurred for balancing of load with demand. Although improvements to grid balancing can be costly, they can lead to long term savings.
In many countries for many types of variable renewable energy, from time to time the government invites companies to tender sealed bids to construct a certain capacity of solar power to connect to certain electricity substations. By accepting the lowest bid the government commits to buy at that price per kWh for a fixed number of years, or up to a certain total amount of power. This provides certainty for investors against highly volatile wholesale electricity prices. However they may still risk exchange rate volatility if they borrowed in foreign currency.
Examples by country
Great Britain
The operator of the British electricity system has said that it will be capable of operating zero-carbon by 2025, whenever there is enough renewable generation, and may be carbon negative by 2033. The company, National Grid Electricity System Operator, states that new products and services will help reduce the overall cost of operating the system.
Germany
In countries with a considerable amount of renewable energy, solar energy causes price drops around noon every day. PV production follows the higher demand during these hours. The images below show two weeks in 2022 in Germany, where renewable energy has a share of over 40%. Prices also drop every night and weekend due to low demand. In hours without PV and wind power, electricity prices rise. This can lead to demand side adjustments. While industry is dependent on the hourly prices, most private households still pay a fixed tariff. With smart meters, private consomers can also be motivated i.e. to load an electric car when enough renewable energy is available and prices are cheap.
Steerable flexibility in electricity production is essential to back up variable energy sources. The German example shows that pumped hydro storage, gas plants and hard coal jump in fast. Lignite varies on a daily basis. Nuclear power and biomass can theoretically adjust to a certain extent. However, in this case incentives still seem not be high enough.
| Technology | Concepts | null |
1427634 | https://en.wikipedia.org/wiki/Rate-determining%20step | Rate-determining step | In chemical kinetics, the overall rate of a reaction is often approximately determined by the slowest step, known as the rate-determining step (RDS or RD-step or r/d step) or rate-limiting step. For a given reaction mechanism, the prediction of the corresponding rate equation (for comparison with the experimental rate law) is often simplified by using this approximation of the rate-determining step.
In principle, the time evolution of the reactant and product concentrations can be determined from the set of simultaneous rate equations for the individual steps of the mechanism, one for each step. However, the analytical solution of these differential equations is not always easy, and in some cases numerical integration may even be required. The hypothesis of a single rate-determining step can greatly simplify the mathematics. In the simplest case the initial step is the slowest, and the overall rate is just the rate of the first step.
Also, the rate equations for mechanisms with a single rate-determining step are usually in a simple mathematical form, whose relation to the mechanism and choice of rate-determining step is clear. The correct rate-determining step can be identified by predicting the rate law for each possible choice and comparing the different predictions with the experimental law, as for the example of and CO below.
The concept of the rate-determining step is very important to the optimization and understanding of many chemical processes such as catalysis and combustion.
Example reaction: + CO
As an example, consider the gas-phase reaction + CO → NO + . If this reaction occurred in a single step, its reaction rate (r) would be proportional to the rate of collisions between and CO molecules: r = k[][CO], where k is the reaction rate constant, and square brackets indicate a molar concentration. Another typical example is the Zel'dovich mechanism.
First step rate-determining
In fact, however, the observed reaction rate is second-order in and zero-order in CO, with rate equation r = k[]2. This suggests that the rate is determined by a step in which two molecules react, with the CO molecule entering at another, faster, step. A possible mechanism in two elementary steps that explains the rate equation is:
+ → NO + (slow step, rate-determining)
+ CO → + (fast step)
In this mechanism the reactive intermediate species is formed in the first step with rate r1 and reacts with CO in the second step with rate r2. However, can also react with NO if the first step occurs in the reverse direction (NO + → 2 ) with rate r−1, where the minus sign indicates the rate of a reverse reaction.
The concentration of a reactive intermediate such as [] remains low and almost constant. It may therefore be estimated by the steady-state approximation, which specifies that the rate at which it is formed equals the (total) rate at which it is consumed. In this example is formed in one step and reacts in two, so that
The statement that the first step is the slow step actually means that the first step in the reverse direction is slower than the second step in the forward direction, so that almost all is consumed by reaction with CO and not with NO. That is, r−1 ≪ r2, so that r1 − r2 ≈ 0. But the overall rate of reaction is the rate of formation of final product (here ), so that r = r2 ≈ r1. That is, the overall rate is determined by the rate of the first step, and (almost) all molecules that react at the first step continue to the fast second step.
Pre-equilibrium: if the second step were rate-determining
The other possible case would be that the second step is slow and rate-determining, meaning that it is slower than the first step in the reverse direction: r2 ≪ r−1. In this hypothesis, r1 − r−1 ≈ 0, so that the first step is (almost) at equilibrium. The overall rate is determined by the second step: r = r2 ≪ r1, as very few molecules that react at the first step continue to the second step, which is much slower. Such a situation in which an intermediate (here ) forms an equilibrium with reactants prior to the rate-determining step is described as a pre-equilibrium For the reaction of and CO, this hypothesis can be rejected, since it implies a rate equation that disagrees with experiment.
+ → NO + (fast step)
+ CO → + (slow step, rate-determining)
If the first step were at equilibrium, then its equilibrium constant expression permits calculation of the concentration of the intermediate in terms of more stable (and more easily measured) reactant and product species:
The overall reaction rate would then be
which disagrees with the experimental rate law given above, and so disproves the hypothesis that the second step is rate-determining for this reaction. However, some other reactions are believed to involve rapid pre-equilibria prior to the rate-determining step, as shown below.
Nucleophilic substitution
Another example is the unimolecular nucleophilic substitution (SN1) reaction in organic chemistry, where it is the first, rate-determining step that is unimolecular. A specific case is the basic hydrolysis of tert-butyl bromide () by aqueous sodium hydroxide. The mechanism has two steps (where R denotes the tert-butyl radical ):
Formation of a carbocation R−Br → + .
Nucleophilic attack by hydroxide ion + → ROH.
This reaction is found to be first-order with r = k[R−Br], which indicates that the first step is slow and determines the rate. The second step with OH− is much faster, so the overall rate is independent of the concentration of OH−.
In contrast, the alkaline hydrolysis of methyl bromide () is a bimolecular nucleophilic substitution (SN2) reaction in a single bimolecular step. Its rate law is second-order: r = k[R−Br][].
Composition of the transition state
A useful rule in the determination of mechanism is that the concentration factors in the rate law indicate the composition and charge of the activated complex or transition state. For the –CO reaction above, the rate depends on []2, so that the activated complex has composition , with 2 entering the reaction before the transition state, and CO reacting after the transition state.
A multistep example is the reaction between oxalic acid and chlorine in aqueous solution: + → 2 + 2 + 2 .
The observed rate law is
which implies an activated complex in which the reactants lose 2 + before the rate-determining step. The formula of the activated complex is + − 2 − + , or (an unknown number of water molecules are added because the possible dependence of the reaction rate on was not studied, since the data were obtained in water solvent at a large and essentially unvarying concentration).
One possible mechanism in which the preliminary steps are assumed to be rapid pre-equilibria occurring prior to the transition state is
+ HOCl + +
+
HOCl + → + + 2
Reaction coordinate diagram
In a multistep reaction, the rate-determining step does not necessarily correspond to the highest Gibbs energy on the reaction coordinate diagram. If there is a reaction intermediate whose energy is lower than the initial reactants, then the activation energy needed to pass through any subsequent transition state depends on the Gibbs energy of that state relative to the lower-energy intermediate. The rate-determining step is then the step with the largest Gibbs energy difference relative either to the starting material or to any previous intermediate on the diagram.
Also, for reaction steps that are not first-order, concentration terms must be considered in choosing the rate-determining step.
Chain reactions
Not all reactions have a single rate-determining step. In particular, the rate of a chain reaction is usually not controlled by any single step.
Diffusion control
In the previous examples the rate determining step was one of the sequential chemical reactions leading to a product. The rate-determining step can also be the transport of reactants to where they can interact and form the product. This case is referred to as diffusion control and, in general, occurs when the formation of product from the activated complex is very rapid and thus the provision of the supply of reactants is rate-determining.
| Physical sciences | Kinetics | Chemistry |
1427763 | https://en.wikipedia.org/wiki/Collision%20theory | Collision theory | Collision theory is a principle of chemistry used to predict the rates of chemical reactions. It states that when suitable particles of the reactant hit each other with the correct orientation, only a certain amount of collisions result in a perceptible or notable change; these successful changes are called successful collisions. The successful collisions must have enough energy, also known as activation energy, at the moment of impact to break the pre-existing bonds and form all new bonds. This results in the products of the reaction. The activation energy is often predicted using the Transition state theory. Increasing the concentration of the reactant brings about more collisions and hence more successful collisions. Increasing the temperature increases the average kinetic energy of the molecules in a solution, increasing the number of collisions that have enough energy. Collision theory was proposed independently by Max Trautz in 1916 and William Lewis in 1918.
When a catalyst is involved in the collision between the reactant molecules, less energy is required for the chemical change to take place, and hence more collisions have sufficient energy for the reaction to occur. The reaction rate therefore increases.
Collision theory is closely related to chemical kinetics.
Collision theory was initially developed for the gas reaction system with no dilution. But most reactions involve solutions, for example, gas reactions in a carrying inert gas, and almost all reactions in solutions. The collision frequency of the solute molecules in these solutions is now controlled by diffusion or Brownian motion of individual molecules. The flux of the diffusive molecules follows Fick's laws of diffusion. For particles in a solution, an example model to calculate the collision frequency and associated coagulation rate is the Smoluchowski coagulation equation proposed by Marian Smoluchowski in a seminal 1916 publication. In this model, Fick's flux at the infinite time limit is used to mimic the particle speed of the collision theory. Jixin Chen proposed a finite-time solution to the diffusion flux in 2022 which significantly changes the estimated collision frequency of two particles in a solution.
Rate equations
The rate for a bimolecular gas-phase reaction, A + B → product, predicted by collision theory is
where:
k is the rate constant in units of (number of molecules)−1⋅s−1⋅m3.
nA is the number density of A in the gas in units of m−3.
nB is the number density of B in the gas in units of m−3. E.g. for a gas mixture with gas A concentration 0.1 mol⋅L−1 and B concentration 0.2 mol⋅L−1, the number of density of A is 0.1×6.02×1023÷10−3 = 6.02×1025 m−3, the number of density of B is 0.2×6.02×1023÷10−3 = 1.2×1026 m−3
Z is the collision frequency in units of m−3⋅s−1.
is the steric factor.
Ea is the activation energy of the reaction, in units of J⋅mol−1.
T is the temperature in units of K.
R is the gas constant in units of J mol−1K−1.
The unit of r(T) can be converted to mol⋅L−1⋅s−1, after divided by (1000×NA), where NA is the Avogadro constant.
For a reaction between A and B, the collision frequency calculated with the hard-sphere model with the unit number of collisions per m3 per second is:
where:
nA is the number density of A in the gas in units of m−3.
nB is the number density of B in the gas in units of m−3. E.g. for a gas mixture with gas A concentration 0.1 mol⋅L−1 and B concentration 0.2 mol⋅L−1, the number of density of A is 0.1×6.02×1023÷10−3 = 6.02×1025 m−3, the number of density of B is 0.2×6.02×1023÷10−3 = 1.2×1026 m−3.
σAB is the reaction cross section (unit m2), the area when two molecules collide with each other, simplified to , where rA the radius of A and rB the radius of B in unit m.
kB is the Boltzmann constant unit J⋅K−1.
T is the absolute temperature (unit K).
μAB is the reduced mass of the reactants A and B, (unit kg).
NA is the Avogadro constant.
[A] is molar concentration of A in unit mol⋅L−1.
[B] is molar concentration of B in unit mol⋅L−1.
Z can be converted to mole collision per liter per second dividing by 1000NA.
If all the units that are related to dimension are converted to dm, i.e. mol⋅dm−3 for [A] and [B], dm2 for σAB, dm2⋅kg⋅s−2⋅K−1 for the Boltzmann constant, then
unit mol⋅dm−3⋅s−1.
Quantitative insights
Derivation
Consider the bimolecular elementary reaction:
A + B → C
In collision theory it is considered that two particles A and B will collide if their nuclei get closer than a certain distance. The area around a molecule A in which it can collide with an approaching B molecule is called the cross section (σAB) of the reaction and is, in simplified terms, the area corresponding to a circle whose radius () is the sum of the radii of both reacting molecules, which are supposed to be spherical.
A moving molecule will therefore sweep a volume per second as it moves, where is the average velocity of the particle. (This solely represents the classical notion of a collision of solid balls. As molecules are quantum-mechanical many-particle systems of electrons and nuclei based upon the Coulomb and exchange interactions, generally they neither obey rotational symmetry nor do they have a box potential. Therefore, more generally the cross section is defined as the reaction probability of a ray of A particles per areal density of B targets, which makes the definition independent from the nature of the interaction between A and B. Consequently, the radius is related to the length scale of their interaction potential.)
From kinetic theory it is known that a molecule of A has an average velocity (different from root mean square velocity) of , where is the Boltzmann constant, and is the mass of the molecule.
The solution of the two-body problem states that two different moving bodies can be treated as one body which has the reduced mass of both and moves with the velocity of the center of mass, so, in this system must be used instead of .
Thus, for a given molecule A, it travels before hitting a molecule B if all B is fixed with no movement, where is the average traveling distance. Since B also moves, the relative velocity can be calculated using the reduced mass of A and B.
Therefore, the total collision frequency, of all A molecules, with all B molecules, is
From Maxwell–Boltzmann distribution it can be deduced that the fraction of collisions with more energy than the activation energy is . Therefore, the rate of a bimolecular reaction for ideal gases will be
in unit number of molecular reactions s−1⋅m−3,
where:
Z is the collision frequency with unit s−1⋅m−3. The z is Z without [A][B].
is the steric factor, which will be discussed in detail in the next section,
Ea is the activation energy (per mole) of the reaction in unit J/mol,
T is the absolute temperature in unit K,
R is the gas constant in unit J/mol/K.
[A] is molar concentration of A in unit mol/L,
[B] is molar concentration of B in unit mol/L.
The product zρ is equivalent to the preexponential factor of the Arrhenius equation.
Validity of the theory and steric factor
Once a theory is formulated, its validity must be tested, that is, compare its predictions with the results of the experiments.
When the expression form of the rate constant is compared with the rate equation for an elementary bimolecular reaction, , it is noticed that
unit M−1⋅s−1 (= dm3⋅mol−1⋅s−1), with all dimension unit dm including kB.
This expression is similar to the Arrhenius equation and gives the first theoretical explanation for the Arrhenius equation on a molecular basis. The weak temperature dependence of the preexponential factor is so small compared to the exponential factor that it cannot be measured experimentally, that is, "it is not feasible to establish, on the basis of temperature studies of the rate constant, whether the predicted T dependence of the preexponential factor is observed experimentally".
Steric factor
If the values of the predicted rate constants are compared with the values of known rate constants, it is noticed that collision theory fails to estimate the constants correctly, and the more complex the molecules are, the more it fails. The reason for this is that particles have been supposed to be spherical and able to react in all directions, which is not true, as the orientation of the collisions is not always proper for the reaction. For example, in the hydrogenation reaction of ethylene the H2 molecule must approach the bonding zone between the atoms, and only a few of all the possible collisions fulfill this requirement.
To alleviate this problem, a new concept must be introduced: the steric factor ρ. It is defined as the ratio between the experimental value and the predicted one (or the ratio between the frequency factor and the collision frequency):
and it is most often less than unity.
Usually, the more complex the reactant molecules, the lower the steric factor. Nevertheless, some reactions exhibit steric factors greater than unity: the harpoon reactions, which involve atoms that exchange electrons, producing ions. The deviation from unity can have different causes: the molecules are not spherical, so different geometries are possible; not all the kinetic energy is delivered into the right spot; the presence of a solvent (when applied to solutions), etc.
{| class="wikitable"
|+ Experimental rate constants compared to the ones predicted by collision theory for gas phase reactions
|-
! Reaction
! A, s−1M−1
! Z, s−1M−1
! Steric factor
|-
| 2ClNO → 2Cl + 2NO || 9.4 || 5.9 || 0.16
|-
| 2ClO → Cl2 + O2 || 6.3 || 2.5 || 2.3
|-
| H2 + C2H4 → C2H6 || 1.24 || 7.3 || 1.7
|-
| Br2 + K → KBr + Br || 1.0 || 2.1 || 4.3
|-
|}
Collision theory can be applied to reactions in solution; in that case, the solvent cage has an effect on the reactant molecules, and several collisions can take place in a single encounter, which leads to predicted preexponential factors being too large. ρ values greater than unity can be attributed to favorable entropic contributions.
{| class="wikitable"
|+ Experimental rate constants compared to the ones predicted by collision theory for reactions in solution
|-
! Reaction
! Solvent
! A, 1011 s−1⋅M−1
! Z, 1011 s−1⋅M−1
! Steric factor
|-
| C2H5Br + OH− || ethanol || 4.30 || 3.86 || 1.11
|-
| C2H5O− + CH3I || ethanol ||2.42 || 1.93 || 1.25
|-
| ClCH2CO2− + OH− || water || 4.55 || 2.86 || 1.59
|-
| C3H6Br2 + I− || methanol || 1.07 || 1.39 || 0.77
|-
| HOCH2CH2Cl + OH− || water ||25.5 || 2.78 || 9.17
|-
| 4-CH3C6H4O− + CH3I || ethanol || 8.49 || 1.99 || 4.27
|-
| CH3(CH2)2Cl + I− || acetone || 0.085 || 1.57|| 0.054
|-
| C5H5N + CH3I || C2H2Cl4 || — || — || 2.0 10
|-
|}
Alternative collision models for diluted solutions
Collision in diluted gas or liquid solution is regulated by diffusion instead of direct collisions, which can be calculated from Fick's laws of diffusion. Theoretical models to calculate the collision frequency in solutions have been proposed by Marian Smoluchowski in a seminal 1916 publication at the infinite time limit, and Jixin Chen in 2022 at a finite-time approximation. A scheme of comparing the rate equations in pure gas and solution is shown in the right figure.
For a diluted solution in the gas or the liquid phase, the collision equation developed for neat gas is not suitable when diffusion takes control of the collision frequency, i.e., the direct collision between the two molecules no longer dominates. For any given molecule A, it has to collide with a lot of solvent molecules, let's say molecule C, before finding the B molecule to react with. Thus the probability of collision should be calculated using the Brownian motion model, which can be approximated to a diffusive flux using various boundary conditions that yield different equations in the Smoluchowski model and the JChen Model.
For the diffusive collision, at the infinite time limit when the molecular flux can be calculated from the Fick's laws of diffusion, in 1916 Smoluchowski derived a collision frequency between molecule A and B in a diluted solution:
where:
is the collision frequency, unit #collisions/s in 1 m3 of solution.
is the radius of the collision cross-section, unit m.
is the relative diffusion constant between A and B, unit m2/s, and .
and are the number concentrations of molecules A and B in the solution respectively, unit #molecule/m3.
or
where:
is in unit mole collisions/s in 1 L of solution.
is the Avogadro constant.
is the radius of the collision cross-section, unit m.
is the relative diffusion constant between A and B, unit m2/s.
and are the molar concentrations of A and B respectively, unit mol/L.
is the diffusive collision rate constant, unit L mol−1 s−1.
There have been a lot of extensions and modifications to the Smoluchowski model since it was proposed in 1916.
In 2022, Chen rationales that because the diffusive flux is evolving over time and the distance between the molecules has a finite value at a given concentration, there should be a critical time to cut off the evolution of the flux that will give a value much larger than the infinite solution Smoluchowski has proposed. So he proposes to use the average time for two molecules to switch places in the solution as the critical cut-off time, i.e., first neighbor visiting time. Although an alternative time could be the mean free path time or the average first passenger time, it overestimates the concentration gradient between the original location of the first passenger to the target. This hypothesis yields a fractal reaction kinetic rate equation of diffusive collision in a diluted solution:
where:
is in unit mole collisions/s in 1 L of solution.
is the Avogadro constant.
is the area of the collision cross-section in unit m2.
is the product of the unitless fractions of reactive surface area on A and B. is the effective adsorption cross-section area.
is the relative diffusion constant between A and B, unit m2/s, and .
and are the molar concentrations of A and B respectively, unit mol/L.
is the diffusive collision rate constant, unit L4/3 mol-4/3 s−1.
| Physical sciences | Kinetics | Chemistry |
1430910 | https://en.wikipedia.org/wiki/Stolon | Stolon | In biology, stolons (from Latin stolō, genitive stolōnis – "branch"), also known as runners, are horizontal connections between parts of an organism. They may be part of the organism, or of its skeleton. Typically, animal stolons are exoskeletons (external skeletons).
In botany
In botany, stolons are plant stems which grow at the soil surface or just below ground that form adventitious roots at the nodes, and new plants from the buds. Stolons are often called runners. Rhizomes, in contrast, are root-like stems that may either grow horizontally at the soil surface or in other orientations underground. Thus, not all horizontal stems are called stolons. Plants with stolons are called stoloniferous.
A stolon is a plant propagation strategy and the complex of individuals formed by a mother plant and all its clones produced from stolons form a single genetic individual, a genet.
Morphology
Stolons may have long or short internodes. The leaves along the stolon are usually very small, but in a few cases such as Stachys sylvatica are normal in size.
Stolons arise from the base of the plant. In strawberries the base is above the soil surface; in many bulb-forming species and plants with rhizomes, the stolons remain underground and form shoots that rise to the surface at the ends or from the nodes. The nodes of the stolons produce roots, often all around the node and hormones produced by the roots cause the stolon to initiate shoots with normal leaves. Typically after the formation of the new plant the stolon dies away in a year or two, while rhizomes persist normally for many years or for the life of the plant, adding more length each year to the ends with active growth. The horizontal growth of stolons results from the interplay of different hormones produced at the growing point and hormones from the main plant, with some studies showing that stolon and rhizome growth are affected by the amount of shady light the plant receives with increased production and branching from plants exposed to mixed shade and sun, while plants in all day sun or all shade produce fewer stolons.
A number of plants have soil-level or above-ground rhizomes, including Iris species and many orchid species.
T. Holm (1929) restricted the term rhizome to a horizontal, usually subterranean, stem that produces roots from its lower surface and green leaves from its apex, developed directly from the plumule of the embryo. He recognized stolons as axillary, subterranean branches that do not bear green leaves but only membranaceous, scale-like ones.
A stolon of grasses is defined as a horizontal stem above or on the soil surface that often roots at the internodes.
Plants with stolons
In some Cyperus species the stolons end with the growth of tubers; the tubers are swollen stolons that form new plants.
Some species of crawling plants can also sprout adventitious roots, but are not considered stoloniferous: a stolon is sprouted from an existing stem and can produce a full individual. Examples of plants that extend through stolons include some species from the genera Riccia, Argentina (silverweed), Cynodon, Fragaria, and Pilosella (Hawkweeds), Zoysia japonica, Ranunculus repens. Plants with long, slender stolons are referred to as sarmentose plants.
Other plants with stolons below the soil surface include many grasses, Ajuga, Mentha, and Stachys. Several species of Irises have stolons attached to their rhizomes, including Iris stolonifera.
Lily-of-the-valley (Convallaria majalis) has rhizomes that grow stolon-like stems called stoloniferous rhizomes or leptomorph rhizomes. A number of plants have stoloniferous rhizomes including Asters. These stolon-like rhizomes are long and thin, with long internodes and indeterminate growth with lateral buds at the node, which mostly remain dormant.
In potatoes, the stolons start to grow within 10 days of plants emerging above ground, with tubers usually beginning to form on the end of the stolons. The tubers are modified stolons that hold food reserves, with a few buds that grow into stems. Since it is not a rhizome it does not generate roots, but the new stem growth that grows to the surface produces roots. | Biology and health sciences | Plant anatomy and morphology: General | Biology |
1431156 | https://en.wikipedia.org/wiki/Agrochemical | Agrochemical | An agrochemical or agrichemical, a contraction of agricultural chemical, is a chemical product used in industrial agriculture. Agrichemical typically refers to biocides (pesticides including insecticides, herbicides, fungicides and nematicides) alongside synthetic fertilizers. It may also include hormones and other chemical growth agents. Though the application of mineral fertilizers and pesticidal chemicals has a long history, the majority of agricultural chemicals were developed from the 19th century, and their use were expanded significantly during the Green Revolution and the late 20th century. Agriculture that uses these chemicals is frequently called conventional agriculture.
Agrochemicals are counted among speciality chemicals. Most agrochemicals are products of the petrochemical industry, where chemicals are derivitatives of fossil fuels. The production and use of agrochemicals contribute substantially to climate change, both through direct emissions during production, and through indirect emissions created from soil ecology problems created by the chemicals.
Agrochemicals, especially when improperly used or released in local environments, have led to a number of public health and environmental issues. Agrochemicals and their production can be significant environmental pollution. Agrochemicals are responsible for significant damage to waterways through runoff, and inproperly stored agrochemicals and agrochemical wastes are responsible for spills, especially during extreme weather events. Following the publication of Rachel Carson's Silent Spring, increased global attention has been paid to these ecological impacts of certain classes of chemicals, both in terms of effects on ecosystems and biodiversity loss. Some farmers choose not to use agrochemicals, with sustainable agriculture approaches such as organic farming or agroecology, avoiding use of pesticides and industrial chemicals, in favor of naturally occurring chemicals.
Categories
Biological action
In most of the cases, agrochemicals refer to pesticides.
Pesticides
Insecticides
Herbicides
Fungicides
Algaecides
Rodenticides
Molluscicides
Nematicides
Fertilisers
Soil conditioners
Liming and acidifying agents
Plant growth regulators
Application method
Fumigants
Penetrant
Application process
Agrochemicals are typically applied to seeds or the field using a variety of different methods.
Seed treatment
Sprayers
Aerial application
Ecology
Many agrochemicals are toxic, and agrichemicals in bulk storage may pose significant environmental and/or health risks, particularly in the event of accidental spills. In many countries, use of agrichemicals is highly regulated. Government-issued permits for purchase and use of approved agrichemicals may be required. Significant penalties can result from misuse, including improper storage resulting in spillage. On farms, proper storage facilities and labeling, emergency clean-up equipment and procedures, and safety equipment and procedures for handling, application and disposal are often subject to mandatory standards and regulations. Usually, the regulations are carried out through the registration process.
For instance, bovine somatotropin, though widely used in the United States, is not approved in Canada and some other jurisdictions as there are concerns for the health of cows using it.
Impacts of pesticides
History
Sumerians from 4500 years ago have said to use insecticides in the form of sulfur compounds. Additionally, the Chinese from about 3200 years ago used mercury and arsenic compounds to control the body lice.
Agrochemicals were introduced to protect crops from pests and enhance crop yields. The most common agrochemicals include pesticides and fertilizers. Chemical fertilizers in the 1960s were responsible for the beginning of the "Green Revolution", where using the same surface of land using intensive irrigation and mineral fertilizers such as nitrogen, phosphorus, and potassium has greatly increased food production. Throughout the 1970s through 1980s, pesticide research continued into producing more selective agrochemicals. Due to the adaptation of pests to these chemicals, more and new agrochemicals were being used, causing side effects in the environment.
Companies
Syngenta was the Chinese owned worldwide leader in agrochemical sales in 2013 at approximately US$10.9 billion, followed by Bayer CropScience, BASF, Dow AgroSciences, Monsanto, and then DuPont with about $3.6 billion. It is still in the worldwide leading position based on sales of year 2019. Based on a statistics by statistica, In 2019, the agrochemical market worldwide was worth approximately $234.2 billion. This is expected to increase to more than $300 billion in 2025.
| Technology | Pest and disease control | null |
1431342 | https://en.wikipedia.org/wiki/Potential%20theory | Potential theory | In mathematics and mathematical physics, potential theory is the study of harmonic functions.
The term "potential theory" was coined in 19th-century physics when it was realized that two fundamental forces of nature known at the time, namely gravity and the electrostatic force, could be modeled using functions called the gravitational potential and electrostatic potential, both of which satisfy Poisson's equation—or in the vacuum, Laplace's equation.
There is considerable overlap between potential theory and the theory of Poisson's equation to the extent that it is impossible to draw a distinction between these two fields. The difference is more one of emphasis than subject matter and rests on the following distinction: potential theory focuses on the properties of the functions as opposed to the properties of the equation. For example, a result about the singularities of harmonic functions would be said to belong to potential theory whilst a result on how the solution depends on the boundary data would be said to belong to the theory of Poisson's equation. This is not a hard and fast distinction, and in practice there is considerable overlap between the two fields, with methods and results from one being used in the other.
Modern potential theory is also intimately connected with probability and the theory of Markov chains. In the continuous case, this is closely related to analytic theory. In the finite state space case, this connection can be introduced by introducing an electrical network on the state space, with resistance between points inversely proportional to transition probabilities and densities proportional to potentials. Even in the finite case, the analogue I-K of the Laplacian in potential theory has its own maximum principle, uniqueness principle, balance principle, and others.
Symmetry
A useful starting point and organizing principle in the study of harmonic functions is a consideration of the symmetries of the Laplace equation. Although it is not a symmetry in the usual sense of the term, we can start with the observation that the Laplace equation is linear. This means that the fundamental object of study in potential theory is a linear space of functions. This observation will prove especially important when we consider function space approaches to the subject in a later section.
As for symmetry in the usual sense of the term, we may start with the theorem that the symmetries of the -dimensional Laplace equation are exactly the conformal symmetries of the -dimensional Euclidean space. This fact has several implications. First of all, one can consider harmonic functions which transform under irreducible representations of the conformal group or of its subgroups (such as the group of rotations or translations). Proceeding in this fashion, one systematically obtains the solutions of the Laplace equation which arise from separation of variables such as spherical harmonic solutions and Fourier series. By taking linear superpositions of these solutions, one can produce large classes of harmonic functions which can be shown to be dense in the space of all harmonic functions under suitable topologies.
Second, one can use conformal symmetry to understand such classical tricks and techniques for generating harmonic functions as the Kelvin transform and the method of images.
Third, one can use conformal transforms to map harmonic functions in one domain to harmonic functions in another domain. The most common instance of such a construction is to relate harmonic functions on a disk to harmonic functions on a half-plane.
Fourth, one can use conformal symmetry to extend harmonic functions to harmonic functions on conformally flat Riemannian manifolds. Perhaps the simplest such extension is to consider a harmonic function defined on the whole of Rn (with the possible exception of a discrete set of singular points) as a harmonic function on the -dimensional sphere. More complicated situations can also happen. For instance, one can obtain a higher-dimensional analog of Riemann surface theory by expressing a multi-valued harmonic function as a single-valued function on a branched cover of Rn or one can regard harmonic functions which are invariant under a discrete subgroup of the conformal group as functions on a multiply connected manifold or orbifold.
Two dimensions
From the fact that the group of conformal transforms is infinite-dimensional in two dimensions and finite-dimensional for more than two dimensions, one can surmise that potential theory in two dimensions is different from potential theory in other dimensions. This is correct and, in fact, when one realizes that any two-dimensional harmonic function is the real part of a complex analytic function, one sees that the subject of two-dimensional potential theory is substantially the same as that of complex analysis. For this reason, when speaking of potential theory, one focuses attention on theorems which hold in three or more dimensions. In this connection, a surprising fact is that many results and concepts originally discovered in complex analysis (such as Schwarz's theorem, Morera's theorem, the Weierstrass-Casorati theorem, Laurent series, and the classification of singularities as removable, poles and essential singularities) generalize to results on harmonic functions in any dimension. By considering which theorems of complex analysis are special cases of theorems of potential theory in any dimension, one can obtain a feel for exactly what is special about complex analysis in two dimensions and what is simply the two-dimensional instance of more general results.
Local behavior
An important topic in potential theory is the study of the local behavior of harmonic functions. Perhaps the most fundamental theorem about local behavior is the regularity theorem for Laplace's equation, which states that harmonic functions are analytic. There are results which describe the local structure of level sets of harmonic functions. There is Bôcher's theorem, which characterizes the behavior of isolated singularities of positive harmonic functions. As alluded to in the last section, one can classify the isolated singularities of harmonic functions as removable singularities, poles, and essential singularities.
Inequalities
A fruitful approach to the study of harmonic functions is the consideration of inequalities they satisfy. Perhaps the most basic such inequality, from which most other inequalities may be derived, is the maximum principle. Another important result is Liouville's theorem, which states the only bounded harmonic functions defined on the whole of Rn are, in fact, constant functions. In addition to these basic inequalities, one has Harnack's inequality, which states that positive harmonic functions on bounded domains are roughly constant.
One important use of these inequalities is to prove convergence of families of harmonic functions or sub-harmonic functions, see Harnack's theorem. These convergence theorems are used to prove the existence of harmonic functions with particular properties.
Spaces of harmonic functions
Since the Laplace equation is linear, the set of harmonic functions defined on a given domain is, in fact, a vector space. By defining suitable norms and/or inner products, one can exhibit sets of harmonic functions which form Hilbert or Banach spaces. In this fashion, one obtains such spaces as the Hardy space, Bloch space, Bergman space and Sobolev space.
| Mathematics | Harmonic analysis | null |
29018709 | https://en.wikipedia.org/wiki/Correlation%20coefficient | Correlation coefficient | A correlation coefficient is a numerical measure of some type of linear correlation, meaning a statistical relationship between two variables. The variables may be two columns of a given data set of observations, often called a sample, or two components of a multivariate random variable with a known distribution.
Several types of correlation coefficient exist, each with their own definition and own range of usability and characteristics. They all assume values in the range from −1 to +1, where ±1 indicates the strongest possible correlation and 0 indicates no correlation. As tools of analysis, correlation coefficients present certain problems, including the propensity of some types to be distorted by outliers and the possibility of incorrectly being used to infer a causal relationship between the variables (for more, see Correlation does not imply causation).
Types
There are several different measures for the degree of correlation in data, depending on the kind of data: principally whether the data is a measurement, ordinal, or categorical.
Pearson
The Pearson product-moment correlation coefficient, also known as , , or Pearson's , is a measure of the strength and direction of the linear relationship between two variables that is defined as the covariance of the variables divided by the product of their standard deviations. This is the best-known and most commonly used type of correlation coefficient. When the term "correlation coefficient" is used without further qualification, it usually refers to the Pearson product-moment correlation coefficient.
Intra-class
Intraclass correlation (ICC) is a descriptive statistic that can be used, when quantitative measurements are made on units that are organized into groups; it describes how strongly units in the same group resemble each other.
Rank
Rank correlation is a measure of the relationship between the rankings of two variables, or two rankings of the same variable:
Spearman's rank correlation coefficient is a measure of how well the relationship between two variables can be described by a monotonic function.
The Kendall tau rank correlation coefficient is a measure of the portion of ranks that match between two data sets.
Goodman and Kruskal's gamma is a measure of the strength of association of the cross tabulated data when both variables are measured at the ordinal level.
Tetrachoric and polychoric
The polychoric correlation coefficient measures association between two ordered-categorical variables. It's technically defined as the estimate of the Pearson correlation coefficient one would obtain if:
The two variables were measured on a continuous scale, instead of as ordered-category variables.
The two continuous variables followed a bivariate normal distribution.
When both variables are dichotomous instead of ordered-categorical, the polychoric correlation coefficient is called the tetrachoric correlation coefficient.
Interpreting correlation coefficient values
The correlation between two variables have different associations that are measured in values such as or . Correlation values range from −1 to +1, where ±1 indicates the strongest possible correlation and 0 indicates no correlation between variables.
| Mathematics | Probability | null |
5210590 | https://en.wikipedia.org/wiki/Lemniscate | Lemniscate | In algebraic geometry, a lemniscate ( or ) is any of several figure-eight or -shaped curves. The word comes from the Latin , meaning "decorated with ribbons", from the Greek (), meaning "ribbon", or which alternatively may refer to the wool from which the ribbons were made.
Curves that have been called a lemniscate include three quartic plane curves: the hippopede or lemniscate of Booth, the lemniscate of Bernoulli, and the lemniscate of Gerono. The hippopede was studied by Proclus (5th century), but the term "lemniscate" was not used until the work of Jacob Bernoulli in the late 17th century.
History and examples
Lemniscate of Booth
The consideration of curves with a figure-eight shape can be traced back to Proclus, a Greek Neoplatonist philosopher and mathematician who lived in the 5th century AD. Proclus considered the cross-sections of a torus by a plane parallel to the axis of the torus. As he observed, for most such sections the cross section consists of either one or two ovals; however, when the plane is tangent to the inner surface of the torus, the cross-section takes on a figure-eight shape, which Proclus called a horse fetter (a device for holding two feet of a horse together), or "hippopede" in Greek. The name "lemniscate of Booth" for this curve dates to its study by the 19th-century mathematician James Booth.
The lemniscate may be defined as an algebraic curve, the zero set of the quartic polynomial when the parameter d is negative (or zero for the special case where the lemniscate becomes a pair of externally tangent circles). For positive values of d one instead obtains the oval of Booth.
Lemniscate of Bernoulli
In 1680, Cassini studied a family of curves, now called the Cassini oval, defined as follows: the locus of all points, the product of whose distances from two fixed points, the curves' foci, is a constant. Under very particular circumstances (when the half-distance between the points is equal to the square root of the constant) this gives rise to a lemniscate.
In 1694, Johann Bernoulli studied the lemniscate case of the Cassini oval, now known as the lemniscate of Bernoulli (shown above), in connection with a problem of "isochrones" that had been posed earlier by Leibniz. Like the hippopede, it is an algebraic curve, the zero set of the polynomial . Bernoulli's brother Jacob Bernoulli also studied the same curve in the same year, and gave it its name, the lemniscate. It may also be defined geometrically as the locus of points whose product of distances from two foci equals the square of half the interfocal distance. It is a special case of the hippopede (lemniscate of Booth), with , and may be formed as a cross-section of a torus whose inner hole and circular cross-sections have the same diameter as each other. The lemniscatic elliptic functions are analogues of trigonometric functions for the lemniscate of Bernoulli, and the lemniscate constants arise in evaluating the arc length of this lemniscate.
Lemniscate of Gerono
Another lemniscate, the lemniscate of Gerono or lemniscate of Huygens, is the zero set of the quartic polynomial . Viviani's curve, a three-dimensional curve formed by intersecting a sphere with a cylinder, also has a figure eight shape, and has the lemniscate of Gerono as its planar projection.
Others
Other figure-eight shaped algebraic curves include
The Devil's curve, a curve defined by the quartic equation in which one connected component has a figure-eight shape,
Watt's curve, a figure-eight shaped curve formed by a mechanical linkage. Watt's curve is the zero set of the degree-six polynomial equation and has the lemniscate of Bernoulli as a special case.
| Mathematics | Two-dimensional space | null |
5217516 | https://en.wikipedia.org/wiki/Clothes%20moth | Clothes moth | Clothes moth or clothing moth is the common name for several species of moth considered to be pests, whose larvae eat animal fibres (hairs), including clothing and other fabrics.
These include:
Tineola bisselliella, the common clothes moth or webbing clothes moth
Tinea pellionella, the case-bearing clothes moth. Obsolete names are: Phalaena (Tinea) pellionella, Phalaena zoolegella, Tinea demiurga, Tinea gerasimovi, and Tinea pelliomella
Trichophaga tapetzella, the carpet moth or tapestry moth
Monopis crocicapitella, pale-backed clothes moth. Particularly destructive of textiles, and found to have increased dramatically in south-west England in 2018.
Niditinea fuscella, the brown-dotted clothes moth
Diet and infestation
The larvae of clothes moths can eat animal fibres which are not removed by other scavengers, and are capable of consuming and digesting keratin materials that make up silk, wool, fur, and hair. This allows clothes moths to attack human-made garments and textiles which include animal fibres, damaging them and leading to the common name of these pests.
Household-wide infestations can stem from a single textile item, such as a garment or rug, with potential targets besides garments including upholstery, toys, or even taxidermied animals. Discarded fibres found around the home can contribute to infestations as well, such as pet sheddings, hair and fur buildup inside vents and air ducts, or birds' nests built inside some part of a house. Basements and attics are commonly the most heavily affected areas. Larvae can also sometimes act as bookworms, chewing through paper (which they cannot digest for nutrition) to reach book bindings or mold colonies for nourishment.
Treatment and control
Various means are used to repel or kill moths. Pheromone traps are also used both to count and to destroy clothes moths, although these only attract certain species of clothes moth so it is possible to have an active clothes moth infestation without any moths being found on the pheromone traps. Care should be taken to correctly distinguish clothes moths from similar-appearing pantry or Indianmeal moths, so that the appropriate pheromone traps are used.
To control an infestation, all items from the area should be removed and cleaned by using a strong-suctioned vacuum, or soapy water and a washing machine to launder items in hot water for 20–30 minutes to kill eggs.
Among other methods, recommendations to protect heritage collections of textiles include checking the undersides of chairs, moving and vacuum-cleaning all furniture once a month and sealing the discarded vacuum cleaner bag, checking and shaking textiles every month, and regularly checking attics and chimneys. If textiles do become infested, adults, eggs and larvae can be killed by freezing garments in sealed bags for a fortnight (14 days).
| Biology and health sciences | Lepidoptera | Animals |
2878626 | https://en.wikipedia.org/wiki/Optical%20computing | Optical computing | Optical computing or photonic computing uses light waves produced by lasers or incoherent sources for data processing, data storage or data communication for computing. For decades, photons have shown promise to enable a higher bandwidth than the electrons used in conventional computers (see optical fibers).
Most research projects focus on replacing current computer components with optical equivalents, resulting in an optical digital computer system processing binary data. This approach appears to offer the best short-term prospects for commercial optical computing, since optical components could be integrated into traditional computers to produce an optical-electronic hybrid. However, optoelectronic devices consume 30% of their energy converting electronic energy into photons and back; this conversion also slows the transmission of messages. All-optical computers eliminate the need for optical-electrical-optical (OEO) conversions, thus reducing electrical power consumption.
Application-specific devices, such as synthetic-aperture radar (SAR) and optical correlators, have been designed to use the principles of optical computing. Correlators can be used, for example, to detect and track objects, and to classify serial time-domain optical data.
Optical components for binary digital computer
The fundamental building block of modern electronic computers is the transistor. To replace electronic components with optical ones, an equivalent optical transistor is required. This is achieved by crystal optics (using materials with a non-linear refractive index). In particular, materials exist where the intensity of incoming light affects the intensity of the light transmitted through the material in a similar manner to the current response of a bipolar transistor. Such an optical transistor can be used to create optical logic gates, which in turn are assembled into the higher level components of the computer's central processing unit (CPU). These will be nonlinear optical crystals used to manipulate light beams into controlling other light beams.
Like any computing system, an optical computing system needs four things to function well:
optical processor
optical data transfer, e.g. fiber-optic cable
optical storage,
optical power source (light source)
Substituting electrical components will need data format conversion from photons to electrons, which will make the system slower.
Controversy
There are some disagreements between researchers about the future capabilities of optical computers; whether or not they may be able to compete with semiconductor-based electronic computers in terms of speed, power consumption, cost, and size is an open question. Critics note that real-world logic systems require "logic-level restoration, cascadability, fan-out and input–output isolation", all of which are currently provided by electronic transistors at low cost, low power, and high speed. For optical logic to be competitive beyond a few niche applications, major breakthroughs in non-linear optical device technology would be required, or perhaps a change in the nature of computing itself.
Misconceptions, challenges, and prospects
A significant challenge to optical computing is that computation is a nonlinear process in which multiple signals must interact. Light, which is an electromagnetic wave, can only interact with another electromagnetic wave in the presence of electrons in a material, and the strength of this interaction is much weaker for electromagnetic waves, such as light, than for the electronic signals in a conventional computer. This may result in the processing elements for an optical computer requiring more power and larger dimensions than those for a conventional electronic computer using transistors.
A further misconception is that since light can travel much faster than the drift velocity of electrons, and at frequencies measured in THz, optical transistors should be capable of extremely high frequencies. However, any electromagnetic wave must obey the transform limit, and therefore the rate at which an optical transistor can respond to a signal is still limited by its spectral bandwidth. In fiber-optic communications, practical limits such as dispersion often constrain channels to bandwidths of tens of GHz, only slightly better than many silicon transistors. Obtaining dramatically faster operation than electronic transistors would therefore require practical methods of transmitting ultrashort pulses down highly dispersive waveguides.
Photonic logic
Photonic logic is the use of photons (light) in logic gates (NOT, AND, OR, NAND, NOR, XOR, XNOR). Switching is obtained using nonlinear optical effects when two or more signals are combined.
Resonators are especially useful in photonic logic, since they allow a build-up of energy from constructive interference, thus enhancing optical nonlinear effects.
Other approaches that have been investigated include photonic logic at a molecular level, using photoluminescent chemicals. In a demonstration, Witlicki et al. performed logical operations using molecules and SERS.
Unconventional approaches
Time delays optical computing
The basic idea is to delay light (or any other signal) in order to perform useful computations. Of interest would be to solve NP-complete problems as those are difficult problems for the conventional computers.
There are two basic properties of light that are actually used in this approach:
The light can be delayed by passing it through an optical fiber of a certain length.
The light can be split into multiple (sub)rays. This property is also essential because we can evaluate multiple solutions in the same time.
When solving a problem with time-delays the following steps must be followed:
The first step is to create a graph-like structure made from optical cables and splitters. Each graph has a start node and a destination node.
The light enters through the start node and traverses the graph until it reaches the destination. It is delayed when passing through arcs and divided inside nodes.
The light is marked when passing through an arc or through a node so that we can easily identify that fact at the destination node.
At the destination node we will wait for a signal (fluctuation in the intensity of the signal) which arrives at a particular moment(s) in time. If there is no signal arriving at that moment, it means that we have no solution for our problem. Otherwise the problem has a solution. Fluctuations can be read with a photodetector and an oscilloscope.
The first problem attacked in this way was the Hamiltonian path problem.
The simplest one is the subset sum problem. An optical device solving an instance with four numbers {a1, a2, a3, a4} is depicted below:
The light will enter in Start node. It will be divided into two (sub)rays of smaller intensity. These two rays will arrive into the second node at moments a1 and 0. Each of them will be divided into two subrays which
will arrive in the third node at moments 0, a1, a2 and a1 + a2. These represents the all subsets of the set {a1, a2}. We expect fluctuations in the intensity of the signal at no more than four different moments. In the destination node we expect fluctuations at no more than 16 different moments (which are all the subsets of the given). If we have a fluctuation in the target moment B, it means that we have a solution of the problem, otherwise there is no subset whose sum of elements equals B. For the practical implementation we cannot have zero-length cables, thus all cables are increased with a small (fixed for all) value k'. In this case the solution is expected at moment B+n×k.
On-Chip Photonic Tensor Cores
With increasing demands on graphical processing unit-based accelerator technologies, in the second decade of the 21st century, there has been a huge emphasis on the use of on-chip integrated optics to create photonics-based processors. The emergence of both deep learning neural networks based on phase modulation, and more recently amplitude modulation using photonic memories have created a new area of photonic technologies for neuromorphic computing, leading to new photonic computing technologies, all on a chip such as the photonic tensor core.
Wavelength-based computing
Wavelength-based computing can be used to solve the 3-SAT problem with n variables, m clauses and with no more than three variables per clause. Each wavelength, contained in a light ray, is considered as possible value-assignments to n variables. The optical device contains prisms and mirrors are used to discriminate proper wavelengths which satisfy the formula.
Computing by xeroxing on transparencies
This approach uses a photocopier and transparent sheets for performing computations. k-SAT problem with n variables, m clauses and at most k variables per clause has been solved in three steps:
Firstly all 2n possible assignments of n variables have been generated by performing n photocopies.
Using at most 2k copies of the truth table, each clause is evaluated at every row of the truth table simultaneously.
The solution is obtained by making a single copy operation of the overlapped transparencies of all m clauses.
Masking optical beams
The travelling salesman problem has been solved by Shaked et al. (2007) by using an optical approach. All possible TSP paths have been generated and stored in a binary matrix which was multiplied with another gray-scale vector containing the distances between cities. The multiplication is performed optically by using an optical correlator.
Optical Fourier co-processors
Many computations, particularly in scientific applications, require frequent use of the 2D discrete Fourier transform (DFT) – for example in solving differential equations describing propagation of waves or transfer of heat. Though modern GPU technologies typically enable high-speed computation of large 2D DFTs, techniques have been developed that can perform continuous Fourier transform optically by utilising the natural Fourier transforming property of lenses. The input is encoded using a liquid crystal spatial light modulator and the result is measured using a conventional CMOS or CCD image sensor. Such optical architectures can offer superior scaling of computational complexity due to the inherently highly interconnected nature of optical propagation, and have been used to solve 2D heat equations.
Ising machines
Physical computers whose design was inspired by the theoretical Ising model are called Ising machines.
Yoshihisa Yamamoto's lab at Stanford pioneered building Ising machines using photons. Initially Yamamoto and his colleagues built an Ising machine using lasers, mirrors, and other optical components commonly found on an optical table.
Later a team at Hewlett Packard Labs developed photonic chip design tools and used them to build an Ising machine on a single chip, integrating 1,052 optical components on that single chip.
Industry
Some additional companies involved with optical computing development include IBM, Microsoft, Procyon Photonics, Lightelligence, Lightmatter, Optalysys, Xanadu Quantum Technologies, QuiX Quantum, ORCA Computing, PsiQuantum, , and TundraSystems Global.
| Technology | Computer architecture concepts | null |
2882046 | https://en.wikipedia.org/wiki/Hass%20avocado | Hass avocado | The Hass avocado is a variety of avocado with dark green, bumpy skin. It was first grown and sold by Southern California mail carrier and amateur horticulturist Rudolph Hass, who also gave it his name.
The Hass avocado is a large-sized fruit weighing 200 to 300 grams (7 to 10 oz). When ripe, the skin becomes a dark purplish-black and yields to gentle pressure. When ripe, the flesh is pale green near the skin and becomes a deeper yellow-green towards the center.
Owing to its taste, size, shelf-life, high growing yield and in some areas, year-round harvesting, the Hass cultivar is the most commercially popular avocado worldwide. In the United States it accounts for more than 80% of the avocado crop and 95% of the California crop, and it is the most widely grown avocado in New Zealand.
History
All commercial, fruit-bearing Hass avocado trees have been grown from grafted seedlings propagated from a single tree that was grown from a seed bought by Rudolph Hass in 1926 from A. R. Rideout of Whittier, California. At the time, Rideout was getting seeds from any source he could find, even restaurant food scraps. The cultivar this seed came from is not known. In 2019, the National Academy of Sciences published a genetic study concluding that the Hass avocado is a cross between Mexican (61%) and Guatemalan (39%) avocado varieties.
In 1926, at his 1.5-acre grove at 430 West Road, La Habra Heights, California, Hass planted three seeds he had bought from Rideout, which yielded one strong seedling. After trying and failing at least twice to graft the seedling with branches from Fuerte avocado trees (the leading commercial cultivar at the time), Hass thought of cutting it down but a professional grafter named Caulkins told him the young tree was sound and strong, so he let it be. When the tree began bearing odd, bumpy fruit, his children liked the taste. As the tree's yields grew bigger, Hass easily sold what his family did not eat to co-workers at the post office. The Hass avocado had one of its first commercial successes at the Model Grocery Store on Colorado Street in Pasadena, California, where chefs working for some of the town's wealthy residents bought the new cultivar's big, nutty-tasting fruit for $1 each, a very high price at the time ().
Hass patented the tree in 1935 and made a contract with Whittier nurseryman Harold Brokaw to grow and sell grafted seedlings propagated from its cuttings, with Brokaw getting 75% of the proceeds. Brokaw then specialized in the Hass and often sold out of grafted seedlings since, unlike the Fuerte, Hass yields are year-round and also more plentiful, with bigger fruit, a longer shelf life and richer flavor owing to higher oil content.
By the early 21st century the US avocado industry took in over $1 billion a year from the heavy-bearing, high quality Hass cultivar, which accounted for around 80% of all avocados grown worldwide.
The mother tree
Owing to later suburban sprawl in Southern California, the mother tree stood for many years in front of a residence in La Habra Heights. The tree died when it was 76 years old and was cut down on 11 September 2002 after a ten-year fight with phytophthora (root rot), which often kills avocado trees. Two plaques by the private residence at 426 West Road mark the spot where it grew. The wood was stored in a tree nursery and from this stock, a nephew of Rudolph Hass, Dick Stewart, made keepsakes, jewelry and other gifts. From 2010 to 2013, in mid-May, and starting again in September 2018, the city of La Habra Heights celebrated the Hass avocado at its Annual La Habra Heights Avocado Festival.
Bearing pattern
Hass avocado trees, like some other cultivars, may only bear well every other year. After a year with low yield, often because of cold, for which the tree does not have much tolerance, yields may be very high the next year. However, the heavy crop can deplete stored carbohydrates, lowering the following season's yield and this can set the tree into a lifelong alternate bearing pattern. Southern California Hass Avocado groves have good soil and drainage, plentiful sunlight and cool gentle winds from the oceans which help the fruit grow. These conditions hold throughout the year, so there are always fresh harvests of Hass avocados in Southern California.
Nutritional value
Raw avocado is 73% water, 15% fat, 9% carbohydrates, and 2% protein (table). As reliable sources are not available for the micronutrient content specifically of Hass avocados, US Department of Agriculture data for a "commercial variety" is used. A 100-gram reference amount supplies of food energy and is rich (20% or higher of the Daily Value, DV) in several B vitamins and vitamin K, with moderate content (10–19% DV) of vitamin C, vitamin E, and potassium (right table, USDA nutrient data). Hass avocados contain phytosterols and carotenoids, including lutein and zeaxanthin.
Avocados have diverse fats. For a typical avocado:
About 75% of an avocado's energy comes from fat, most of which (67% of total fat) is monounsaturated fat as oleic acid.
Other predominant fats include palmitic acid and linoleic acid.
The saturated fat content amounts to 14% of the total fat.
Typical total fat composition is roughly: 1% ω-3, 14% ω-6, 71% ω-9 (65% oleic and 6% palmitoleic), and 14% saturated fat (palmitic acid).
| Biology and health sciences | Other culinary fruits | Plants |
24458151 | https://en.wikipedia.org/wiki/Planetary%20boundaries | Planetary boundaries | Planetary boundaries are a framework to describe limits to the impacts of human activities on the Earth system. Beyond these limits, the environment may not be able to self-regulate anymore. This would mean the Earth system would leave the period of stability of the Holocene, in which human society developed. The framework is based on scientific evidence that human actions, especially those of industrialized societies since the Industrial Revolution, have become the main driver of global environmental change. According to the framework, "transgressing one or more planetary boundaries may be deleterious or even catastrophic due to the risk of crossing thresholds that will trigger non-linear, abrupt environmental change within continental-scale to planetary-scale systems."
The normative component of the framework is that human societies have been able to thrive under the comparatively stable climatic and ecological conditions of the Holocene. To the extent that these Earth system process boundaries have not been crossed, they mark the "safe zone" for human societies on the planet. Proponents of the planetary boundary framework propose returning to this environmental and climatic system; as opposed to human science and technology deliberately creating a more beneficial climate. The concept doesn't address how humans have massively altered ecological conditions to better suit themselves. The climatic and ecological Holocene this framework considers as a "safe zone" doesn't involve massive industrial farming. So this framework begs a reassessment of how to feed modern populations.
The concept has since become influential in the international community (e.g. United Nations Conference on Sustainable Development), including governments at all levels, international organizations, civil society and the scientific community. The framework consists of nine global change processes. In 2009, according to Rockström and others, three boundaries were already crossed (biodiversity loss, climate change and nitrogen cycle), while others were in imminent danger of being crossed.
In 2015, several of the scientists in the original group published an update, bringing in new co-authors and new model-based analysis. According to this update, four of the boundaries were crossed: climate change, loss of biosphere integrity, land-system change, altered biogeochemical cycles (phosphorus and nitrogen). The scientists also changed the name of the boundary "Loss of biodiversity" to "Change in biosphere integrity" to emphasize that not only the number of species but also the functioning of the biosphere as a whole is important for Earth system stability. Similarly, the "Chemical pollution" boundary was renamed to "Introduction of novel entities", widening the scope to consider different kinds of human-generated materials that disrupt Earth system processes.
In 2022, based on the available literature, the introduction of novel entities was concluded to be the 5th transgressed planetary boundary. Freshwater change was concluded to be the 6th transgressed planetary boundary in 2023.
Framework overview and principles
The basic idea of the Planetary Boundaries framework is that maintaining the observed resilience of the Earth system in the Holocene is a precondition for humanity's pursuit of long-term social and economic development. The Planetary Boundaries framework contributes to an understanding of global sustainability because it brings a planetary scale and a long timeframe into focus.
The framework described nine "planetary life support systems" essential for maintaining a "desired Holocene state", and attempted to quantify how far seven of these systems had been pushed already. Boundaries were defined to help define a "safe space for human development", which was an improvement on approaches aiming at minimizing human impacts on the planet.
The framework is based on scientific evidence that human actions, especially those of industrialized societies since the Industrial Revolution, have become the main driver of global environmental change. According to the framework, "transgressing one or more planetary boundaries may be deleterious or even catastrophic due to the risk of crossing thresholds that will trigger non-linear, abrupt environmental change within continental-scale to planetary-scale systems." The framework consists of nine global change processes. In 2009, two boundaries were already crossed, while others were in imminent danger of being crossed. Later estimates indicated that three of these boundaries—climate change, biodiversity loss, and the biogeochemical flow boundary—appear to have been crossed.
The scientists outlined how breaching the boundaries increases the threat of functional disruption, even collapse, in Earth's biophysical systems in ways that could be catastrophic for human wellbeing. While they highlighted scientific uncertainty, they indicated that breaching boundaries could "trigger feedbacks that may result in crossing thresholds that drastically reduce the ability to return within safe levels". The boundaries were "rough, first estimates only, surrounded by large uncertainties and knowledge gaps" which interact in complex ways that are not yet well understood.
The planetary boundaries framework lays the groundwork for a shifting approach to governance and management, away from the essentially sectoral analyses of limits to growth aimed at minimizing negative externalities, toward the estimation of the safe space for human development. Planetary boundaries demarcate, as it were, the "planetary playing field" for humanity if major human-induced environmental change on a global scale is to be avoided.
Authors
The authors of this framework was a group of Earth System and environmental scientists in 2009 led by Johan Rockström from the Stockholm Resilience Centre and Will Steffen from the Australian National University. They collaborated with 26 leading academics, including Nobel laureate Paul Crutzen, Goddard Institute for Space Studies climate scientist James Hansen, oceanographer Katherine Richardson, geographer Diana Liverman and the German Chancellor's chief climate adviser Hans Joachim Schellnhuber.
Most of the contributing scientists were involved in strategy-setting for the Earth System Science Partnership, the precursor to the international global change research network Future Earth. The group wanted to define a "safe operating space for humanity" for the wider scientific community, as a precondition for sustainable development.
Nine boundaries
Thresholds and tipping points
The 2009 study identified nine planetary boundaries and, drawing on current scientific understanding, the researchers proposed quantifications for seven of them. These are:
climate change (CO2 concentration in the atmosphere < 350 ppm and/or a maximum change of +1 W/m2 in radiative forcing);
ocean acidification (mean surface seawater saturation state with respect to aragonite ≥ 80% of pre-industrial levels);
stratospheric ozone depletion (less than 5% reduction in total atmospheric O3 from a pre-industrial level of 290 Dobson Units);
biogeochemical flows in the nitrogen (N) cycle (limit industrial and agricultural fixation of N2 to 35 Tg N/yr) and phosphorus (P) cycle (annual P inflow to oceans not to exceed 10 times the natural background weathering of P);
global freshwater use (< 4000 km3/yr of consumptive use of runoff resources);
land system change (< 15% of the ice-free land surface under cropland);
the erosion of biosphere integrity (an annual rate of loss of biological diversity of < 10 extinctions per million species).
chemical pollution (introduction of novel entities in the environment).
For one process in the planetary boundaries framework, the scientists have not specified a global boundary quantification:
atmospheric aerosol loading;
The quantification of individual planetary boundaries is based on the observed dynamics of the interacting Earth system processes included in the framework. The control variables were chosen because together they provide an effective way to track the human-caused shift away from Holocene conditions.
For some of Earth's dynamic processes, historic data display clear thresholds between comparatively stable conditions. For example, past ice-ages show that during peak glacial conditions, the atmospheric concentration of CO2 was ~180-200 ppm. In interglacial periods (including the Holocene), CO2 concentration has fluctuated around 280 ppm. To know what past climate conditions were like with an atmosphere with over 350 ppm CO2, scientists need to look back about 3 million years. The paleo record of climatic, ecological and biogeochemical changes shows that the Earth system has experienced tipping points, when a very small increment for a control variable (like CO2) triggers a larger, possibly catastrophic, change in the response variable (global warming) through feedbacks in the natural Earth System itself.
For several of the processes in the planetary boundaries framework, it is difficult to locate individual points that mark the threshold shift away from Holocene-like conditions. This is because the Earth system is complex and the scientific evidence base is still partial and fragmented. Instead, the planetary boundaries framework identifies many Earth system thresholds at multiple scales that will be influenced by increases in the control variables. Examples include shifts in monsoon behavior linked to the aerosol loading and freshwater use planetary boundaries.
"Safe operating spaces"
The planetary boundaries framework proposes a range of values for its control variables. This range is supposed to span the threshold between a 'safe operating space' where Holocene-like dynamics can be maintained and a highly uncertain, poorly predictable world where Earth system changes likely increase risks to societies. The boundary is defined as the lower end of that range. If the boundaries are persistently crossed, the world goes further into a danger zone.
It is difficult to restore a 'safe operating space' for humanity that is described by the planetary boundary concept. Even if past biophysical changes could be mitigated, the predominant paradigms of social and economic development appear largely indifferent to the looming possibilities of large scale environmental disasters triggered by human actions. Legal boundaries can help keep human activities in check, but are only as effective as the political will to make and enforce them.
Interaction among boundaries
Understanding the Earth system is fundamentally about understanding interactions among environmental change processes. The planetary boundaries are defined with reference to dynamic conditions of the Earth system, but scientific discussions about how different planetary boundaries relate to each other are often philosophically and analytically muddled. Clearer definitions of the basic concepts and terms might help give clarity.
There are many many interactions among the processes in the planetary boundaries framework. While these interactions can create both stabilizing and destabilizing feedbacks in the Earth system, the authors suggested that a transgressed planetary boundary will reduce the safe operating space for other processes in the framework rather than expand it from the proposed boundary levels. They give the example that the land use boundary could "shift downward" if the freshwater boundary is breached, causing lands to become arid and unavailable for agriculture. At a regional level, water resources may decline in Asia if deforestation continues in the Amazon. That way of framing the interactions shifts from the framework's biophysical definition of boundaries based on Holocene-like conditions to an anthropocentric definition (demand for agricultural land). Despite this conceptual slippage, considerations of known Earth system interactions across scales suggest the need for "extreme caution in approaching or transgressing any individual planetary boundaries."
Another example has to do with coral reefs and marine ecosystems: In 2009, researchers showed that, since 1990, calcification in the reefs of the Great Barrier that they examined decreased at a rate unprecedented over the last 400 years (14% in less than 20 years). Their evidence suggests that the increasing temperature stress and the declining ocean saturation state of aragonite is making it difficult for reef corals to deposit calcium carbonate. Multiple stressors, such as increased nutrient loads and fishing pressure, moves corals into less desirable ecosystem states. Ocean acidification will significantly change the distribution and abundance of a whole range of marine life, particularly species "that build skeletons, shells, and tests of biogenic calcium carbonate. Increasing temperatures, surface UV radiation levels and ocean acidity all stress marine biota, and the combination of these stresses may well cause perturbations in the abundance and diversity of marine biological systems that go well beyond the effects of a single stressor acting alone."
Proposed new or expanded boundaries since 2012
In 2012, Steven Running suggested a tenth boundary, the annual net global primary production of all terrestrial plants, as an easily determinable measure integrating many variables that will give "a clear signal about the health of ecosystems".
In 2015, a second paper was published in Science to update the Planetary Boundaries concept. The update concluded four boundaries had now been transgressed: climate, biodiversity, land use and biogeochemical cycles. The 2015 paper emphasized interactions of the nine boundaries and identified climate change and loss of biodiversity integrity as 'core boundaries' of central importance to the framework because the interactions of climate and the biosphere are what scientifically defines Earth system conditions.
In 2017, some authors argued that marine systems are underrepresented in the framework. Their proposed remedy was to include the seabed as a component of the earth surface change boundary. They also wrote that the framework should account for "changes in vertical mixing and ocean circulation patterns".
Subsequent work on planetary boundaries begins to relate these thresholds at the regional scale.
Debate and further research per boundary
Climate change
A 2018 study calls into question the adequacy of efforts to limit warming to 2 °C above pre-industrial temperatures, as set out in the Paris Agreement. The scientists raise the possibility that even if greenhouse gas emissions are substantially reduced to limit warming to 2 °C, that might exceed the "threshold" at which self-reinforcing climate feedbacks add additional warming until the climate system stabilizes in a hothouse climate state. This would make parts of the world uninhabitable for people, raise sea levels by up to , and raise temperatures by to levels that are higher than any interglacial period in the past 1.2 million years.
Change in biosphere integrity
According to the biologist Cristián Samper, a "boundary that expresses the probability of families of species disappearing over time would better reflect our potential impacts on the future of life on Earth." The biodiversity boundary has also been criticized for framing biodiversity solely in terms of the extinction rate. The global extinction rate has been highly variable over the Earth's history, and thus using it as the only biodiversity variable can be of limited usefulness.
Nitrogen and phosphorus
The biogeochemist William Schlesinger thinks waiting until we near some suggested limit for nitrogen deposition and other pollutions will just permit us to continue to a point where it is too late. He says the boundary suggested for phosphorus is not sustainable, and would exhaust the known phosphorus reserves in less than 200 years.
The ocean chemist Peter Brewer queries whether it is "truly useful to create a list of environmental limits without serious plans for how they may be achieved ... they may become just another stick to beat citizens with. Disruption of the global nitrogen cycle is one clear example: it is likely that a large fraction of people on Earth would not be alive today without the artificial production of fertilizer. How can such ethical and economic issues be matched with a simple call to set limits? ... food is not optional."
Peak phosphorus is a concept to describe the point in time at which the maximum global phosphorus production rate is reached. Phosphorus is a scarce finite resource on earth and means of production other than mining are unavailable because of its non-gaseous environmental cycle. According to some researchers, Earth's phosphorus reserves are expected to be completely depleted in 50–100 years and peak phosphorus to be reached by approximately 2030. However, recent evidence shows that if phosphorus applications to soil are matched to the agronomic optimum for crop yield, it would take >500 years to exhaust currently econimically viable phosphorus reserves.
Ocean acidification
Surface ocean acidity is clearly interconnected with the climate change boundaries, since the concentration of carbon dioxide in the atmosphere is also the underlying control variable for the ocean acidification boundary.
The ocean chemist Peter Brewer thinks "ocean acidification has impacts other than simple changes in pH, and these may need boundaries too."
Land-system change
Across the planet, forests, wetlands and other vegetation types are being converted to agricultural and other land uses, impacting freshwater, carbon and other cycles, and reducing biodiversity. In the year 2015 the boundary was defined as 75% of forests rested intact, including 85% of boreal forests, 50% of temperate forests and 85% of tropical forests. The boundary is crossed because only 62% of forests rested intact as of the year 2015.
The boundary for land use has been criticized as follows: "The boundary of 15 per cent land-use change is, in practice, a premature policy guideline that dilutes the authors' overall scientific proposition. Instead, the authors might want to consider a limit on soil degradation or soil loss. This would be a more valid and useful indicator of the state of terrestrial health."
Freshwater
The freshwater cycle is another boundary significantly affected by climate change. Overexploitation of freshwater occurs if a water resource is mined or extracted at a rate that exceeds the recharge rate. Water pollution and saltwater intrusion can also turn much of the world's underground water and lakes into finite resources with "peak water" usage debates similar to oil.
The hydrologist David Molden stated in 2009 that planetary boundaries are a welcome new approach in the "limits to growth" debate but said "a global limit on water consumption is necessary, but the suggested planetary boundary of 4,000 cubic kilometres per year is too generous."
Green and blue water
A study concludes that the 'Freshwater use' boundary should be renamed to the 'Freshwater change', composed of "green" and "blue" water components. 'Green water' refers to disturbances of terrestrial precipitation, evaporation and soil moisture. Water scarcity can have substantial effects in agriculture. When measuring and projecting water scarcity in agriculture for climate change scenarios, both "green water" and "blue water" are of relevance.
In April 2022, scientists proposed and preliminarily evaluated 'green water' in the water cycle as a likely transgressed planetary boundary, as measured by root-zone soil moisture deviation from Holocene variability.
Ozone depletion
The stratospheric ozone layer protectively filters ultraviolet radiation (UV) from the Sun, which would otherwise damage biological systems. The actions taken after the Montreal Protocol appeared to be keeping the planet within a safe boundary.
The Nobel laureate in chemistry, Mario Molina, says "five per cent is a reasonable limit for acceptable ozone depletion, but it doesn't represent a tipping point".
Atmospheric aerosols
Worldwide each year, aerosol particles result in about 800,000 premature deaths from air pollution. Aerosol loading is sufficiently important to be included among the planetary boundaries, but it is not yet clear whether an appropriate safe threshold measure can be identified.
Novel entities (chemical pollution)
Some chemicals, such as persistent organic pollutants, heavy metals and radionuclides, have potentially irreversible additive and synergic effects on biological organisms, reducing fertility and resulting in permanent genetic damage. Sublethal uptakes are drastically reducing marine bird and mammal populations. This boundary seems important, although it is hard to quantify. In 2019, it was suggested that novel entities could include genetically modified organisms, pesticides and even artificial intelligence.
A Bayesian emulator for persistent organic pollutants has been developed which can potentially be used to quantify the boundaries for chemical pollution. To date, critical exposure levels of polychlorinated biphenyls (PCBs) above which mass mortality events of marine mammals are likely to occur, have been proposed as a chemical pollution planetary boundary.
There are at least 350,000 artificial chemicals in the world. They are coming from "plastics, pesticides, industrial chemicals, chemicals in consumer products, antibiotics and other pharmaceuticals". They have mostly "negative effects on planetary health". Their production increased 50 times since 1950 and is expected to increase 3 times more by 2050. Plastic alone contain more than 10,000 chemicals and create large problems. The researchers are calling for limit on chemical production and shift to circular economy, meaning to products that can be reused and recycled.
In January 2022 a group of scientists concluded that this planetary boundary is already exceeded, which puts in risk the stability of the Earth system. They integrated the literature information on how production and release of a number of novel entities, including plastics and hazardous chemicals, have rapidly increased in the last decades with significant impact on the planetary processes.
In August 2022, scientists concluded that the (overall transgressed) boundary is a placeholder for multiple different boundaries for NEs that may emerge, reporting that PFAS pollution is one such new boundary. They show that levels of these so-called "forever chemicals" in rainwater are ubiquitously, and often greatly, above guideline safe levels worldwide. There are some moves to restrict and replace their use.
Related concepts
Planetary integrity
Planetary integrity is also called earth's life-support systems or ecological integrity. Scholars have pointed out that planetary integrity "needs to be maintained for long-term sustainability". The current biodiversity loss is threatening ecological integrity on a global scale. The term integrity refers to ecological health in this context. The concept of planetary integrity is interlinked within the concept of planetary boundaries.
An expert Panel on Ecological Integrity in 1998 has defined ecological integrity as follows: "Ecosystems have integrity when they have their native components (plants, animals and other organisms) and processes (such as growth and reproduction) intact."
The Sustainable Development Goals might be able to act as a steering mechanism to address the current loss of planetary integrity. There are many negative human impacts on the environment that are causing a reduction in planetary integrity.
The "Limits to Growth" (1972) and Gaia theory
The idea that there are limits to the burden placed upon our planet by human activities has been around for a long time. The Planetary Boundaries framework acknowledges the influence of the 1972 study, The Limits to Growth, that presented a model in which exponential growth in world population, industrialization, pollution, food production, and resources depletion outstrip the ability of technology to increase resources availability. Subsequently, the report was widely dismissed, particularly by economists and business people, and it has often been claimed that history has proved the projections to be incorrect. In 2008, Graham Turner from the Commonwealth Scientific and Industrial Research Organisation (CSIRO) published "A comparison of The Limits to Growth with thirty years of reality". The Limits to Growth has been widely discussed, both by critics of the modelling approach and its conclusions and by analysts who argue that the insight that societies do not live in an unlimited world and that historical data since the 1970s support the report's findings. The Limits to Growth approach explores how the socio-technical dynamics of the world economy may limit humanity's opportunities and introduce risks of collapse. In contrast, the Planetary Boundaries framework focuses on the biophysical dynamics of the Earth system.
Our Common Future was published in 1987 by United Nations' World Commission on Environment and Development. It tried to recapture the spirit of the Stockholm Conference. Its aim was to interlock the concepts of development and environment for future political discussions. It introduced the famous definition for sustainable development: "Development that meets the needs of the present without compromising the ability of future generations to meet their own needs."
Another key idea influencing the Planetary Boundaries framework is the Gaia theory or hypothesis. In the 1970s, James Lovelock and microbiologist Lynn Margulis presented the idea that all organisms and their inorganic surroundings on Earth are integrated into a single self-regulating system. The system has the ability to react to perturbations or deviations, much like a living organism adjusts its regulation mechanisms to accommodate environmental changes such as temperature (homeostasis). Nevertheless, this capacity has limits. For instance, when a living organism is subjected to a temperature that is lower or higher than its living range, it can perish because its regulating mechanism cannot make the necessary adjustments. Similarly the Earth may not be able to react to large deviations in critical parameters. In Lovelock's book The Revenge of Gaia, he suggests that the destruction of rainforests and biodiversity, compounded with global warming resulting from the increase of greenhouse gases made by humans, could shift feedbacks in the Earth system away from a self-regulating balance to a positive (intensifying) feedback loop.
Anthropocene
Scientists have affirmed that the planet has entered a new epoch, the Anthropocene. In the Anthropocene, humans have become the main agents of not only change to the Earth System but also the driver of Earth System rupture, disruption of the Earth System's ability to be resilient and recover from that change, potentially ultimately threatening planetary habitability. The previous geological epoch, the Holocene began about 10,000 years ago. It is the current interglacial period, and was a relatively stable environment of the Earth. There have been natural environmental fluctuations during the Holocene, but the key atmospheric and biogeochemical parameters have remained within relatively narrow bounds. This stability has allowed societies to thrive worldwide, developing agriculture, large-scale settlements and complex networks of trade.
According to Rockström et al., we "have now become so dependent on those investments for our way of life, and how we have organized society, technologies, and economies around them, that we must take the range within which Earth System processes varied in the Holocene as a scientific reference point for a desirable planetary state."
Various biophysical processes that are important in maintaining the resilience of the Earth system are also undergoing large and rapid change because of human actions. For example, since the advent of the Anthropocene, the rate at which species are going extinct has increased over 100 times, and humans are now the driving force altering global river flows as well as water vapor flows from the land surface. Continuing perturbation of Earth system processes by human activities raises the possibility that further pressure could be destabilizing, leading to non-linear, abrupt, large-scale or irreversible environmental change responses by the Earth system within continental- to planetary-scale systems.
Reception and debate
The 2009 report was presented to the General Assembly of the Club of Rome in Amsterdam. An edited summary of the report was published as the featured article in a special 2009 edition of Nature alongside invited critical commentary from leading academics like Nobel laureate Mario J. Molina and biologist Cristián Samper.
Development studies scholars have been critical of aspects of the framework and constraints that its adoption could place on the Global South. Proposals to conserve a certain proportion of Earth's remaining forests can be seen as rewarding the countries such as those in Europe that have already economically benefitted from exhausting their forests and converting land for agriculture. In contrast, countries that have yet to industrialize are asked to make sacrifices for global environmental damage they may have had little role in creating.
The biogeochemist William Schlesinger queries whether thresholds are a good idea for pollutions at all. He thinks waiting until we near some suggested limit will just permit us to continue to a point where it is too late. "Management based on thresholds, although attractive in its simplicity, allows pernicious, slow and diffuse degradation to persist nearly indefinitely."
In a global empirical study, researchers investigated how students of environmental and sustainability studies in 35 countries assessed the planetary boundaries. It was found that there are substantial global differences in the perception of planetary boundaries.
Subsequent developments
The "safe and just space" doughnut
National environmental footprints
Several studies have assessed environmental footprints of nations based on planetary boundaries: for Portugal, Sweden, Switzerland, the Netherlands, the European Union, India, many of Belt and Road Initiative countries as well as for the world's most important economies. While the metrics and allocation approaches applied varied, there is a converging outcome that resource use of wealthier nations – if extrapolated to world population – is not compatible with planetary boundaries.
Boundaries related to agriculture and food consumption
Human activities related to agriculture and nutrition globally contribute to the transgression of four out of nine planetary boundaries. Surplus nutrient flows (N, P) into aquatic and terrestrial ecosystems are of highest importance, followed by excessive land-system change and biodiversity loss. Whereas in the case of biodiversity loss, P cycle and land-system change, the transgression is in the zone of uncertainty—indicating an increasing risk (yellow circle in the figure), the N boundary related to agriculture is more than 200% transgressed—indicating a high risk (red marked circle in the figure). Here, nutrition includes food processing and trade as well as food consumption (preparation of food in households and gastronomy). Consumption-related environmental impacts are not quantified at the global level for the planetary boundaries of freshwater use, atmospheric aerosol loading (air pollution) and stratospheric ozone depletion.
Individual and collective allowances
Approaches based on a general framework of ecological limits include (transferable) personal carbon allowances and "legislated" national greenhouse gas emissions limits. Consumers would have freedom in their (informed) choice within (the collective) boundaries.
Usage at international policy level
United Nations
The United Nations secretary general Ban Ki-moon endorsed the concept of planetary boundaries on 16 March 2012, when he presented the key points of the report of his High Level Panel on Global Sustainability to an informal plenary of the UN General Assembly. Ban stated: "The Panel's vision is to eradicate poverty and reduce inequality, to make growth inclusive and production and consumption more sustainable, while combating climate change and respecting a range of other planetary boundaries." The concept was incorporated into the so-called "zero draft" of the outcome of the United Nations Conference on Sustainable Development to be convened in Rio de Janeiro 20–22 June 2012. However, the use of the concept was subsequently withdrawn from the text of the conference, "partly due to concerns from some poorer countries that its adoption could lead to the sidelining of poverty reduction and economic development. It is also, say observers, because the idea is simply too new to be officially adopted, and needed to be challenged, weathered and chewed over to test its robustness before standing a chance of being internationally accepted at UN negotiations."
In 2011, at their second meeting, the High-level Panel on Global Sustainability of the United Nations had incorporated the concept of planetary boundaries into their framework, stating that their goal was: "To eradicate poverty and reduce inequality, make growth inclusive, and production and consumption more sustainable while combating climate change and respecting the range of other planetary boundaries."
Elsewhere in their proceedings, panel members have expressed reservations about the political effectiveness of using the concept of "planetary boundaries": "Planetary boundaries are still an evolving concept that should be used with caution [...] The planetary boundaries question can be divisive as it can be perceived as a tool of the "North" to tell the "South" not to follow the resource intensive and environmentally destructive development pathway that rich countries took themselves... This language is unacceptable to most of the developing countries as they fear that an emphasis on boundaries would place unacceptable brakes on poor countries."
However, the concept is routinely used in the proceedings of the United Nations, and in the UN Daily News. For example, the United Nations Environment Programme (UNEP) Executive Director Achim Steiner states that the challenge of agriculture is to "feed a growing global population without pushing humanity's footprint beyond planetary boundaries." The UNEP Yearbook 2010 also repeated Rockström's message, conceptually linking it with ecosystem management and environmental governance indicators.
In their 2012 report entitled "Resilient People, Resilient Planet: A future worth choosing", The High-level Panel on Global Sustainability called for bold global efforts, "including launching a major global scientific initiative, to strengthen the interface between science and policy. We must define, through science, what scientists refer to as "planetary boundaries", "environmental thresholds" and "tipping points"".
European Commission
The planetary boundaries concept is also used in proceedings by the European Commission, and was referred to in the European Environment Agency synthesis report The European environment – state and outlook 2010.
| Physical sciences | Earth science basics: General | Earth science |
24462957 | https://en.wikipedia.org/wiki/Hazard | Hazard | A hazard is a potential source of harm. Substances, events, or circumstances can constitute hazards when their nature would potentially allow them to cause damage to health, life, property, or any other interest of value. The probability of that harm being realized in a specific incident, combined with the magnitude of potential harm, make up its risk. This term is often used synonymously in colloquial speech.
Hazards can be classified in several ways which are not mutually exclusive. They can be classified by causing actor (for example, natural or anthropogenic), by physical nature (e.g. biological or chemical) or by type of damage (e.g., health hazard or environmental hazard). Examples of natural disasters with highly harmful impacts on a society are floods, droughts, earthquakes, tropical cyclones, lightning strikes, volcanic activity and wildfires. Technological and anthropogenic hazards include, for example, structural collapses, transport accidents, accidental or intentional explosions, and release of toxic materials.
The term climate hazard is used in the context of climate change. These are hazards that stem from climate-related events and can be associated with global warming, such as wildfires, floods, droughts, sea level rise. Climate hazards can combine with other hazards and result in compound event losses (see also loss and damage). For example, the climate hazard of heat can combine with the hazard of poor air quality. Or the climate hazard flooding can combine with poor water quality.
In physics terms, common theme across many forms of hazards is the presence of energy that can cause damage, as it can happen with chemical energy, mechanical energy or thermal energy. This damage can affect different valuable interests, and the severity of the associated risk varies.
Definition
A hazard is defined as "the potential occurrence of a natural or human-induced physical event or trend that may cause loss of life, injury, or other health impacts, as well as damage and loss to property, infrastructure, livelihoods, service provision, ecosystems and environmental resources."
A hazard only exists if there is a pathway to exposure. As an example, the center of the Earth consists of molten material at very high temperatures which would be a severe hazard if contact was made with the core. However, there is no feasible way of making contact with the core, therefore the center of the Earth currently poses no hazard.
The frequency and severity of hazards are important aspects for risk management. Hazards may also be assessed in relation to the impact that they have.
In defining hazard Keith Smith argues that what may be defined as the hazard is only a hazard if there is the presence of humans to make it a hazard. In this regard, human sensitivity to environmental hazards is a combination of both physical exposure (natural and/or technological events at a location related to their statistical variability) and human vulnerability (about social and economic tolerance of the same location).
Relationship with other terms
Disaster
An example of the distinction between a natural hazard and a disaster is that an earthquake is the hazard which caused the 1906 San Francisco earthquake disaster.
A natural disaster is the highly harmful impact on a society or community following a natural hazard event. The term "disaster" itself is defined as follows: "Disasters are serious disruptions to the functioning of a community that exceed its capacity to cope using its own resources. Disasters can be caused by natural, man-made and technological hazards, as well as various factors that influence the exposure and vulnerability of a community."
The US Federal Emergency Management Agency (FEMA) explains the relationship between natural disasters and natural hazards as follows: "Natural hazards and natural disasters are related but are not the same. A natural hazard is the threat of an event that will likely have a negative impact. A natural disaster is the negative impact following an actual occurrence of natural hazard in the event that it significantly harms a community.
Disaster can take various forms, including hurricane, volcano, tsunami, earthquake, drought, famine, plague, disease, rail crash, car crash, tornado, deforestation, flooding, toxic release, and spills (oil, chemicals).
A disaster hazard is an extreme geophysical event that is capable of causing a disaster. 'Extreme' in this case means a substantial variation in either the positive or the negative direction from the normal trend; flood disasters can result from exceptionally high precipitation and river discharge, and drought is caused by exceptionally low values. The fundamental determinants of hazard and the risk of such hazards occurring is timing, location, magnitude and frequency. For example, magnitudes of earthquakes are measured on the Richter scale from 1 to 10, whereby each increment of 1 indicates a tenfold increase in severity. The magnitude-frequency rule states that over a significant period of time many small events and a few large ones will occur. Hurricanes and typhoons on the other hand occur between 5 degrees and 25 degrees north and south of the equator, tending to be seasonal phenomena that are thus largely recurrent in time and predictable in location due to the specific climate variables necessary for their formation.
Risk and vulnerability
The terms hazard and risk are often used interchangeably. However, in terms of risk assessment, these are two very distinct terms. A hazard is an agent that can cause harm or damage to humans, property, or the environment. Risk is the probability that exposure to a hazard will lead to a negative consequence, or more simply, a hazard poses no risk if there is no exposure to that hazard.
Risk is a combination of hazard, exposure and vulnerability. For example in terms of water security: examples of hazards are droughts, floods and decline in water quality. Bad infrastructure and bad governance lead to high exposure to risk.
Risk can be defined as the likelihood or probability of a given hazard of a given level causing a particular level of loss of damage. The elements of risk are populations, communities, the built environment, the natural environment, economic activities and services which are under threat of disaster in a given area.
Another definition of risk is "the probable frequency and probable magnitude of future losses". This definition also focuses on the probability of future loss whereby the degree of vulnerability to hazard represents the level of risk on a particular population or environment. The threats posed by a hazard are:
Hazards to people – death, injury, disease and stress
Hazards to goods – property damage and economic loss
Hazards to environment –loss of flora and fauna, pollution and loss of amenity
Classifications
Hazards can be classified in several ways. These categories are not mutually exclusive which means that one hazard can fall into several categories. For example, water pollution with toxic chemicals is an anthropogenic hazard as well as an environmental hazard.
One of the classification methods is by specifying the origin of the hazard. One key concept in identifying a hazard is the presence of stored energy that, when released, can cause damage. The stored energy can occur in many forms: chemical, mechanical, thermal, radioactive, electrical, etc.
The United Nations Office for Disaster Risk Reduction (UNDRR) explains that "each hazard is characterized by its location, intensity or magnitude, frequency and probability".
A distinction can also be made between rapid-onset natural hazards, technological hazards, and social hazards, which are described as being of sudden occurrence and relatively short duration, and the consequences of longer-term environmental degradation such as desertification and drought.
Hazards may be grouped according to their characteristics. These factors are related to geophysical events, which are not process specific:
Areal extent of damage zone
Intensity of impact at a point
Duration of impact at a point
Rate of onset of the event
Predictability of the event
By causing actor
Natural hazard
Damage to valuable human interests can occur due to phenomena and processes of the natural environment. Natural disasters such as earthquakes, floods, volcanoes and tsunami have threatened people, society, the natural environment, and the built environment, particularly more vulnerable people, throughout history, and in some cases, on a day-to-day basis. According to the Red Cross, each year 130,000 people are killed, 90,000 are injured and 140 million are affected by unique events known as natural disasters.
Potentially dangerous phenomena which are natural or predominantly natural (for example, exceptions are intentional floods) can be classified in these categories:
Meteorological and hydrological hazards, e.g. lightning, storm, flood, sandstorm, fog, rogue wave, tsunami, snow, cold wave, heat wave
Geological hazards
Earthquake
Volcanism, which can cause a wide range of events, such as lava flow and ash fall
Surface or near-surface events, especially erosion and mass wasting (e.g. landslide)
Extraterrestrial hazards, e.g. solar storm, impact event
Natural hazards can be influenced by human actions in different ways and to varying degrees, e.g. land-use change, drainage and construction. Humans play a central role in the existence of natural hazards because "it is only when people and their possessions get in the way of natural processes that hazard exists".
A natural hazard can be considered as a geophysical event when it occurs in extremes and a human factor is involved that may present a risk. There may be an acceptable variation of magnitude which can vary from the estimated normal or average range with upper and lower limits or thresholds. In these extremes, the natural occurrence may become an event that presents a risk to the environment or people. For example, above-average wind speeds resulting in a tropical depression or hurricane according to intensity measures on the Saffir–Simpson scale will provide an extreme natural event that may be considered a hazard.
Seismic hazard
Tsunamis can be caused by geophysical hazards, such as in the 2004 Indian Ocean earthquake and tsunami.
Although generally a natural phenomenon, earthquakes can sometimes be induced by human interventions, such as injection wells, large underground nuclear explosions, excavation of mines, or reservoirs.
Volcanic hazard
Anthropogenic hazard
Anthropogenic hazards, or human-induced hazards, are "induced entirely or predominantly by human activities and choices". These can be societal, technological or environmental hazards.
Technological hazard
Technological hazards are created by the possibility of failure associated with human technology (including emerging technologies), which can also impact the economy, health and national security.
For example, technological hazards can arise from the following events:
Transport accidents: traffic collisions, navigation accidents, rail accidents, aviation accidents
Nuclear materials-related accidents: nuclear plant accidents, nuclear weapon accidents and other nuclear accidents
Chemical accidents and accidental explosions
Mining accidents
Space accidents
Technological failures (including those that affect information and communications technology, such as cybersecurity threats) can cause disruptions to the energy industry (e.g. power outages), telecommunications (e.g. Internet outage), healthcare, banking, transportation, food supply, water supply and other important services.
Structural failures or construction accidents
A mechanical hazard is any hazard involving a machine or industrial process. Motor vehicles, aircraft, and air bags pose mechanical hazards. Compressed gases or liquids can also be considered a mechanical hazard. Hazard identification of new machines and/or industrial processes occurs at various stages in the design of the new machine or process. These hazard identification studies focus mainly on deviations from the intended use or design and the harm that may occur as a result of these deviations. These studies are regulated by various agencies such as the Occupational Safety and Health Administration and the National Highway Traffic Safety Administration.
Engineering hazards occur when human structures fail (e.g. building or structural collapse, bridge failures, dam failures) or the materials used in their construction prove to be hazardous.
Societal hazard
Societal hazards can arise from civil disorders, explosive remnants of war, violence, crowd accidents, financial crises, etc. However, the United Nations Office for Disaster Risk Reduction (UNDRR) Hazard Definition & Classification Review (Sendai Framework 2015 - 2030) specifically excludes armed conflict from the anthropogenic hazard category, as these hazards are already recognised under international humanitarian law.
Waste disposal
In managing waste many hazardous materials are put in the domestic and commercial waste stream. In part this is because modern technological living uses certain toxic or poisonous materials in the electronics and chemical industries. Which, when they are in use or transported, are usually safely contained or encapsulated and packaged to avoid any exposure. In the waste stream, the waste products exterior or encapsulation breaks or degrades and there is a release and exposure to hazardous materials into the environment, for people working in the waste disposal industry, those living around sites used for waste disposal or landfill and the general environment surrounding such sites.
Socionatural hazard
There are different ways to group hazards by origin. The definition by UNDRR states: "Hazards may be natural, anthropogenic or socionatural in origin." The socionatural hazards are those that are "associated with a combination of natural and anthropogenic factors, including environmental degradation and climate change".
Climate hazard
The term climate hazard or climatic hazard is used in the context of climate change, for example in the IPCC Sixth Assessment Report. These are hazards that stem from climate-related events such as wildfires, floods, droughts, sea level rise.
Climate hazards in the context of water include: Increased temperatures, changes in rainfall patterns between the wet and dry season (increased rainfall variability) and sea level rise. The reason why increasing temperatures is listed here as a climate hazard is because "warming temperatures may result in higher evapotranspiration, in turn leading to drier soils".
Waterborne diseases are also connected to climate hazards.
Climate hazards can combine with other hazards and result in compound event losses (see also loss and damage). For example, the climate hazard of heat can combine with the hazard of poor air quality. Or the climate hazard flooding can combine with poor water quality.
Climate scientists have pointed out that climate hazards affect different groups of people differently, depending on their climate change vulnerability: There are "factors that make people and groups vulnerable (e.g., poverty, uneven power structures, disadvantage and discrimination due to, for example, social location and the intersectionality or the overlapping and compounding risks from ethnicity or racial discrimination, gender, age, or disability, etc.)".
By physical nature
Biological hazard
Biological hazards, also known as biohazards, originate in biological processes of living organisms and pose threats to the health of humans, the security of property, or the environment.
Biological hazards include pathogenic microorganisms, such as viruses and bacteria, epidemics, pandemics, parasites, pests, animal attacks, venomous animals, biological toxins and foodborne illnesses.
For example, naturally occurring bacteria such as Escherichia coli and Salmonella are well known pathogens, and a variety of measures have been taken to limit human exposure to these microorganisms through food safety, good personal hygiene, and education. The potential for new biological hazards also exists through the discovery of new microorganisms and the development of new genetically modified (GM) organisms. The use of new GM organisms is regulated by various governmental agencies. The US Environmental Protection Agency (EPA) controls GM plants that produce or resist pesticides (i.e. Bt corn and Roundup ready crops). The US Food and Drug Administration (FDA) regulates GM plants that will be used as food or for medicinal purposes.
Biological hazards can include medical waste or samples of a microorganism, virus or toxin (from a biological source) that can affect health. Many biological hazards are associated with food, including certain viruses, parasites, fungi, bacteria, and plant and seafood toxins. Pathogenic Campylobacter and Salmonella are common foodborne biological hazards. The hazards from these bacteria can be avoided through risk mitigation steps such as proper handling, storing, and cooking of food. Diseases can be enhanced by human factors such as poor sanitation or by processes such as urbanization.
Chemical hazard
A chemical can be considered a hazard if by its intrinsic properties it can cause harm or danger to humans, property, or the environment. Health hazards associated with chemicals are dependent on the dose or amount of the chemical. For example, iodine in the form of potassium iodate is used to produce iodised salt. When applied at a rate of 20 mg of potassium iodate per 1000 mg of table salt, the chemical is beneficial in preventing goitre, while iodine intakes of 1200–9500 mg in one dose has been known to cause death. Some chemicals have a cumulative biological effect, while others are metabolically eliminated over time. Other chemical hazards may depend on concentration or total quantity for their effects.
Some harmful chemicals occur naturally in certain geological formations, such as arsenic. Other chemicals include products with commercial uses, such as agricultural and industrial chemicals, as well as products developed for home use.
A variety of chemical hazards have been identified. However, every year companies produce more new chemicals to fill new needs or to take the place of older, less effective chemicals. Laws, such as the Federal Food, Drug, and Cosmetic Act and the Toxic Substances Control Act in the US, require protection of human health and the environment for any new chemical introduced. In the US, the EPA regulates new chemicals that may have environmental impacts (i.e., pesticides or chemicals released during a manufacturing process), while the FDA regulates new chemicals used in foods or as drugs. The potential hazards of these chemicals can be identified by performing a variety of tests before the authorization of usage. The number of tests required and the extent to which the chemicals are tested varies, depending on the desired usage of the chemical. Chemicals designed as new drugs must undergo more rigorous tests than those used as pesticides.
Pesticides, which are normally used to control unwanted insects and plants, may cause a variety of negative effects on non-target organisms. DDT can build up, or bioaccumulate, in birds, resulting in thinner-than-normal eggshells, which can break in the nest. The organochlorine pesticide dieldrin has been linked to Parkinson's disease. Corrosive chemicals like sulfuric acid, which is found in car batteries and research laboratories, can cause severe skin burns. Many other chemicals used in industrial and laboratory settings can cause respiratory, digestive, or nervous system problems if they are inhaled, ingested, or absorbed through the skin. The negative effects of other chemicals, such as alcohol and nicotine, have been well documented.
Organohalogens are a family of synthetic organic molecules which all contain atoms of one of the halogens. Such materials include PCBs, Dioxins, DDT, Freon and many others. Although considered harmless when first produced, many of these compounds are now known to have profound physiological effects on many organisms including man. Many are also fat soluble and become concentrated through the food chain.
Radioactive or electromagnetic hazard
Radioactive materials produce ionizing radiation which may be very harmful to living organisms. Damage from even a short exposure to radioactivity may have long-term adverse health consequences.
Thermal or fire hazard
Fire hazard
Threats to fire safety are commonly referred to as fire hazards. A fire hazard may include a situation that increases the likelihood of a fire or may impede escape in the event a fire occurs.
Casualties resulting from fires, regardless of their source or initial cause, can be aggravated by inadequate emergency preparedness. Such hazards as a lack of accessible emergency exits, poorly marked escape routes, or improperly maintained fire extinguishers or sprinkler systems may result in many more deaths and injuries than might occur with such protections.
Kinetic hazard
Kinetic energy is involved in hazards associated with noise, falling, or vibration.
By type of damage
Health hazard
Hazards that would affect the health of exposed persons, usually having an acute or chronic illness as the consequence. Fatality would not normally be an immediate consequence. Health hazards may cause measurable changes in the body which are generally indicated by the development of signs and symptoms in the exposed persons, or non-measurable, subjective symptoms.
Ergonomic hazard
Ergonomic hazards are physical conditions that may pose a risk of injury to the musculoskeletal system, such as the muscles or ligaments of the lower back, tendons or nerves of the hands/wrists, or bones surrounding the knees. Ergonomic hazards include things such as awkward or extreme postures, whole-body or hand/arm vibration, poorly designed tools, equipment, or workstations, repetitive motion, and poor lighting. Ergonomic hazards occur in both occupational and non-occupational settings such as in workshops, building sites, offices, home, school, or public spaces and facilities.
Occupational hazard
Psychosocial hazard
Psychological or psychosocial hazards are hazards that affect the psychological well-being of people, including their ability to participate in a work environment among other people. Psychosocial hazards are related to the way work is designed, organized, and managed, as well as the economic and social contexts of work, and are associated with psychiatric, psychological, and/or physical injury or illness. Linked to psychosocial risks are issues such as occupational stress and workplace violence, which are recognized internationally as major challenges to occupational health and safety.
Environmental hazard
Property
Cultural property
Cultural property can be damaged, lost or destroyed by different events or processes, including war, vandalism, theft, looting, transport accident, water leak, human error, natural disaster, fire, pests, pollution and progressive deterioration.
By status
Hazards are sometimes classified into three modes or statuses:
Dormant—The situation environment is currently affected. For instance, a hillside may be unstable, with the potential for a landslide, but there is nothing below or on the hillside that could be affected.
Armed—People, property, or environment are in potential harm's way.
Active—A harmful incident involving the hazard has actually occurred. Often this is referred to not as an "active hazard" but as an accident, emergency, incident, or disaster.
Analysis and management
A range of methodologies are used to assess hazards and to manage them:
(HACCP)
(HAZOP)
Hazard symbol
| Biology and health sciences | Health and fitness: General | Health |
33087778 | https://en.wikipedia.org/wiki/Pressurized%20heavy-water%20reactor | Pressurized heavy-water reactor | A pressurized heavy-water reactor (PHWR) is a nuclear reactor that uses heavy water (deuterium oxide D2O) as its coolant and neutron moderator. PHWRs frequently use natural uranium as fuel, but sometimes also use very low enriched uranium. The heavy water coolant is kept under pressure to avoid boiling, allowing it to reach higher temperature (mostly) without forming steam bubbles, exactly as for a pressurized water reactor (PWR). While heavy water is very expensive to isolate from ordinary water (often referred to as light water in contrast to heavy water), its low absorption of neutrons greatly increases the neutron economy of the reactor, avoiding the need for enriched fuel. The high cost of the heavy water is offset by the lowered cost of using natural uranium and/or alternative fuel cycles. As of the beginning of 2001, 31 PHWRs were in operation, having a total capacity of 16.5 GW(e), representing roughly 7.76% by number and 4.7% by generating capacity of all current operating reactors. CANDU and IPHWR are the most common type of reactors in the PHWR family.
Purpose of using heavy water
The key to maintaining a nuclear chain reaction within a nuclear reactor is to use, on average, exactly one of the neutrons released from each nuclear fission event to stimulate another nuclear fission event (in another fissionable nucleus). With careful design of the reactor's geometry, and careful control of the substances present so as to influence the reactivity, a self-sustaining chain reaction or "criticality" can be achieved and maintained.
Natural uranium consists of a mixture of various isotopes, primarily 238U and a much smaller amount (about 0.72% by weight) of 235U. 238U can only be fissioned by neutrons that are relatively energetic, about 1 MeV or above. No amount of 238U can be made "critical" since it will tend to parasitically absorb more neutrons than it releases by the fission process. 235U, on the other hand, can support a self-sustained chain reaction, but due to the low natural abundance of 235U, natural uranium cannot achieve criticality by itself.
The trick to achieving criticality using only natural or low enriched uranium, for which there is no "bare" critical mass, is to slow down the emitted neutrons (without absorbing them) to the point where enough of them may cause further nuclear fission in the small amount of 235U which is available. (238U which is the bulk of natural uranium is also fissionable with fast neutrons.) This requires the use of a neutron moderator, which absorbs virtually all of the neutrons' kinetic energy, slowing them down to the point that they reach thermal equilibrium with surrounding material. It has been found beneficial to the neutron economy to physically separate the neutron energy moderation process from the uranium fuel itself, as 238U has a high probability of absorbing neutrons with intermediate kinetic energy levels, a reaction known as "resonance" absorption. This is a fundamental reason for designing reactors with separate solid fuel segments, surrounded by the moderator, rather than any geometry that would give a homogeneous mix of fuel and moderator.
Water makes an excellent moderator; the ordinary hydrogen or protium atoms in the water molecules are very close in mass to a single neutron, and so their collisions result in a very efficient transfer of momentum, similar conceptually to the collision of two billiard balls. However, as well as being a good moderator, ordinary water is also quite effective at absorbing neutrons. And so using ordinary water as a moderator will easily absorb so many neutrons that too few are left to sustain a chain reaction with the small isolated 235U nuclei in the fuel, thus precluding criticality in natural uranium. Because of this, a light-water reactor will require that the 235U isotope be concentrated in its uranium fuel, as enriched uranium, generally between 3% and 5% 235U by weight (the by-product from this process enrichment process is known as depleted uranium, and so consisting mainly of 238U, chemically pure). The degree of enrichment needed to achieve criticality with a light-water moderator depends on the exact geometry and other design parameters of the reactor.
One complication of this approach is the need for uranium enrichment facilities, which are generally expensive to build and operate. They also present a nuclear proliferation concern; the same systems used to enrich the 235U can also be used to produce much more "pure" weapons-grade material (90% or more 235U), suitable for producing a nuclear weapon. This is not a trivial exercise by any means, but feasible enough that enrichment facilities present a significant nuclear proliferation risk.
An alternative solution to the problem is to use a moderator that does not absorb neutrons as readily as water. In this case potentially all of the neutrons being released can be moderated and used in reactions with the 235U, in which case there is enough 235U in natural uranium to sustain criticality. One such moderator is heavy water, or deuterium-oxide. Although it reacts dynamically with the neutrons in a fashion similar to light water (albeit with less energy transfer on average, given that heavy hydrogen, or deuterium, is about twice the mass of hydrogen), it already has the extra neutron that light water would normally tend to absorb.
Advantages and disadvantages
Advantages
The use of heavy water as the moderator is the key to the PHWR (pressurized heavy water reactor) system, enabling the use of natural uranium as the fuel (in the form of ceramic UO2), which means that it can be operated without expensive uranium enrichment facilities. The mechanical arrangement of the PHWR, which places most of the moderator at lower temperatures, is particularly efficient because the resulting thermal neutrons have lower energies (neutron temperature after successive passes through a moderator roughly equals the temperature of the moderator) than in traditional designs, where the moderator normally is much hotter. The neutron cross section for fission is higher in the lower the neutron temperature is, and thus lower temperatures in the moderator make successful interaction between neutrons and fissile material more likely. These features mean that a PHWR can use natural uranium and other fuels, and does so more efficiently than light water reactors (LWRs). CANDU type PHWRs are claimed to be able to handle fuels including reprocessed uranium or even spent nuclear fuel from "conventional" light water reactors as well as MOX fuel and there is ongoing research into the ability of CANDU type reactors to operate exclusively on such fuels in a commercial setting. (More on that in the article on the CANDU reactor itself)
Disadvantages
Pressurised heavy-water reactors do have some drawbacks. Heavy water generally costs hundreds of dollars per kilogram, though this is a trade-off against reduced fuel costs. The reduced energy content of natural uranium as compared to enriched uranium necessitates more frequent replacement of fuel; this is normally accomplished by use of an on-power refuelling system. The increased rate of fuel movement through the reactor also results in higher volumes of spent fuel than in LWRs employing enriched uranium. Since unenriched uranium fuel accumulates a lower density of fission products than enriched uranium fuel, however, it generates less heat, allowing more compact storage. While deuterium has a lower neutron capture cross section than protium, this value isn't zero and thus part of the heavy water moderator will inevitably be converted to tritiated water. While tritium, a radioactive isotope of hydrogen, is also produced as a fission product in minute quantities in other reactors, tritium can more easily escape to the environment if it is also present in the cooling water, which is the case in those PHWRs which use heavy water both as moderator and as coolant. Some CANDU reactors separate out the tritium from their heavy water inventory at regular intervals and sell it at a profit, however.
While with typical CANDU derived fuel bundles, the reactor design has a slightly positive Void coefficient of reactivity, the Argentina designed CARA fuel bundles used in Atucha I, are capable of the preferred negative coefficient.
Nuclear proliferation
While prior to India's development of nuclear weapons (see below), the ability to use natural uranium (and thus forego the need for uranium enrichment which is a dual use technology) was seen as hindering nuclear proliferation, this opinion has changed drastically in light of the ability of several countries to build atomic bombs out of plutonium, which can easily be produced in heavy water reactors. Heavy-water reactors may thus pose a greater risk of nuclear proliferation versus comparable light-water reactors due to the low neutron absorption properties of heavy water, discovered in 1937 by Hans von Halban and Otto Frisch. Occasionally, when an atom of 238U is exposed to neutron radiation, its nucleus will capture a neutron, changing it to 239U. The 239U then rapidly undergoes two β− decays — both emitting an electron and an antineutrino, the first one transmuting the 239U into 239Np, and the second one transmuting the 239Np into 239Pu. Although this process takes place with natural uranium using other moderators such as ultra-pure graphite or beryllium, heavy water is by far the best. The Manhattan Project ultimately used graphite moderated reactors to produce plutonium, while the German wartime nuclear project wrongfully dismissed graphite as a suitable moderator due to overlooking impurities and thus made unsuccessful attempts using heavy water (which they correctly identified as an excellent moderator). The Soviet nuclear program likewise used graphite as a moderator and ultimately developed the graphite moderated RBMK as a reactor capable of producing both large amounts of electric power and weapons grade plutonium without the need for heavy water or - at least according to initial design specifications - uranium enrichment.
239Pu is a fissile material suitable for use in nuclear weapons. As a result, if the fuel of a heavy-water reactor is changed frequently, significant amounts of weapons-grade plutonium can be chemically extracted from the irradiated natural uranium fuel by nuclear reprocessing.
In addition, the use of heavy water as a moderator results in the production of small amounts of tritium when the deuterium nuclei in the heavy water absorb neutrons, a very inefficient reaction. Tritium is essential for the production of boosted fission weapons, which in turn enable the easier production of thermonuclear weapons, including neutron bombs. This process is currently expected to provide (at least partially) tritium for ITER.
The proliferation risk of heavy-water reactors was demonstrated when India produced the plutonium for Operation Smiling Buddha, its first nuclear weapon test, by extraction from the spent fuel of a heavy-water research reactor known as the CIRUS reactor.
| Technology | Power generation | null |
33092462 | https://en.wikipedia.org/wiki/Cirrostratus%20nebulosus | Cirrostratus nebulosus | Cirrostratus nebulosus is a type of high-level cirrostratus cloud. The name cirrostratus nebulosus is derived from Latin, the adjective nebulosus meaning "full of vapor, foggy, cloudy, dark". Cirrostratus nebulosus is one of the two most common forms that cirrostratus often takes, with the other being cirrostratus fibratus. The nebulosus species is featureless and uniform, while the fibratus species has a fibrous appearance. Cirrostratus nebulosus are formed by gently rising air. The cloud is often hard to see unless the sun shines through it at the correct angle, forming a halo. While usually very light, the cloud may also be very dense, and the exact appearance of the cloud can vary from one formation to another. In the winter, precipitation often follows behind these clouds; however, they are not a precipitation-producing cloud.
| Physical sciences | Clouds | Earth science |
33094374 | https://en.wikipedia.org/wiki/Telecommunications | Telecommunications | Telecommunication, often used in its plural form or abbreviated as telecom, is the transmission of information over a distance using electronic means, typically through cables, radio waves, or other communication technologies. These means of transmission may be divided into communication channels for multiplexing, allowing for a single medium to transmit several concurrent communication sessions. Long-distance technologies invented during the 20th and 21st centuries generally use electric power, and include the telegraph, telephone, television, and radio.
Early telecommunication networks used metal wires as the medium for transmitting signals. These networks were used for telegraphy and telephony for many decades. In the first decade of the 20th century, a revolution in wireless communication began with breakthroughs including those made in radio communications by Guglielmo Marconi, who won the 1909 Nobel Prize in Physics. Other early pioneers in electrical and electronic telecommunications include co-inventors of the telegraph Charles Wheatstone and Samuel Morse, numerous inventors and developers of the telephone including Antonio Meucci, Philipp Reis, Elisha Gray and Alexander Graham Bell, inventors of radio Edwin Armstrong and Lee de Forest, as well as inventors of television like Vladimir K. Zworykin, John Logie Baird and Philo Farnsworth.
Since the 1960s, the proliferation of digital technologies has meant that voice communications have gradually been supplemented by data. The physical limitations of metallic media prompted the development of optical fibre. The Internet, a technology independent of any given medium, has provided global access to services for individual users and further reduced location and time limitations on communications.
Definition
At the 1932 Plenipotentiary Telegraph Conference and the International Radiotelegraph Conference in Madrid, the two organizations merged to form the International Telecommunication Union (ITU). They defined telecommunication as "any telegraphic or telephonic communication of signs, signals, writing, facsimiles and sounds of any kind, by wire, wireless or other systems or processes of electric signaling or visual signaling (semaphores)."
The definition was later reconfirmed, according to Article 1.3 of the ITU Radio Regulations, which defined it as "Any transmission, emission or reception of signs, signals, writings, images and sounds or intelligence of any nature by wire, radio, optical, or other electromagnetic systems".
As such, slow communications technologies like postal mail and pneumatic tubes are excluded from the telecommunication's definition.
Telecommunication is a compound noun of the Greek prefix tele- (τῆλε), meaning distant, far off, or afar, and the Latin verb communicare, meaning to share. Its modern use is adapted from the French, because its written use was recorded in 1904 by the French engineer and novelist Édouard Estaunié. Communication was first used as an English word in the late 14th century. It comes from Old French comunicacion (14c., Modern French communication), from Latin communicationem (nominative communication), noun of action from past participle stem of communicare, "to share, divide out; communicate, impart, inform; join, unite, participate in," literally, "to make common", from communis".
History
Many transmission media have been used for long-distance communication throughout history, from smoke signals, beacons, semaphore telegraphs, signal flags, and optical heliographs to wires and empty space made to carry electromagnetic signals.
Before the electrical and electronic era
Homing pigeons have been used throughout history by different cultures. Pigeon post had Persian roots and was later used by the Romans to aid their military. Frontinus claimed Julius Caesar used pigeons as messengers in his conquest of Gaul. The Greeks also conveyed the names of the victors at the Olympic Games to various cities using homing pigeons. In the early 19th century, the Dutch government used the system in Java and Sumatra. And in 1849, Paul Julius Reuter started a pigeon service to fly stock prices between Aachen and Brussels, a service that operated for a year until the gap in the telegraph link was closed.
In the Middle Ages, chains of beacons were commonly used on hilltops as a means of relaying a signal. Beacon chains suffered the drawback that they could only pass a single bit of information, so the meaning of the message such as "the enemy has been sighted" had to be agreed upon in advance. One notable instance of their use was during the Spanish Armada, when a beacon chain relayed a signal from Plymouth to London.
In 1792, Claude Chappe, a French engineer, built the first fixed visual telegraphy system (or semaphore line) between Lille and Paris. However semaphore suffered from the need for skilled operators and expensive towers at intervals of ten to thirty kilometres (six to nineteen miles). As a result of competition from the electrical telegraph, the last commercial line was abandoned in 1880.
Besides beacons and pigeons, long-distance communication used sounds like coded drumbeats, the blowing of horns, and whistles.
Telegraph and telephone
On July 25, 1837, the first commercial electrical telegraph was demonstrated by English inventor Sir William Fothergill Cooke and English scientist Sir Charles Wheatstone. Both inventors viewed their device as "an improvement to the [existing] electromagnetic telegraph" and not as a new device.
Samuel Morse independently developed a version of the electrical telegraph that he unsuccessfully demonstrated on September 2, 1837. His code was an important advance over Wheatstone's signaling method. The first transatlantic telegraph cable was successfully completed on July 27, 1866, allowing transatlantic telecommunication for the first time.
The conventional telephone was patented by Alexander Bell in 1876. Elisha Gray also filed a caveat for it in 1876. Gray abandoned his caveat and because he did not contest Bell's priority, the examiner approved Bell's patent on March 3, 1876. Gray had filed his caveat for the variable resistance telephone, but Bell was the first to document the idea and test it in a telephone. Antonio Meucci invented a device that allowed the electrical transmission of voice over a line nearly 30 years before in 1849, but his device was of little practical value because it relied on the electrophonic effect requiring users to place the receiver in their mouths to hear. The first commercial telephone services were set up by the Bell Telephone Company in 1878 and 1879 on both sides of the Atlantic in the cities of New Haven and London.
Radio and television
In 1894, Italian inventor Guglielmo Marconi began developing a wireless communication using the then-newly discovered phenomenon of radio waves, demonstrating, by 1901, that they could be transmitted across the Atlantic Ocean. This was the start of wireless telegraphy by radio. On 17 December 1902, a transmission from the Marconi station in Glace Bay, Nova Scotia, Canada, became the world's first radio message to cross the Atlantic from North America. In 1904, a commercial service was established to transmit nightly news summaries to subscribing ships, which incorporated them into their onboard newspapers.
World War I accelerated the development of radio for military communications. After the war, commercial radio AM broadcasting began in the 1920s and became an important mass medium for entertainment and news. World War II again accelerated the development of radio for the wartime purposes of aircraft and land communication, radio navigation, and radar. Development of stereo FM broadcasting of radio began in the 1930s in the United States and the 1940s in the United Kingdom, displacing AM as the dominant commercial standard in the 1970s.
On March 25, 1925, John Logie Baird demonstrated the transmission of moving pictures at the London department store Selfridges. Baird's device relied upon the Nipkow disk by Paul Nipkow and thus became known as the mechanical television. It formed the basis of experimental broadcasts done by the British Broadcasting Corporation beginning on 30 September 1929. However, for most of the 20th century, televisions depended on the cathode ray tube invented by Karl Ferdinand Braun. The first version of such a television to show promise was produced by Philo Farnsworth and demonstrated to his family on 7 September 1927. After World War II, interrupted experiments resumed and television became an important home entertainment broadcast medium.
Thermionic valves
The type of device known as a thermionic tube or thermionic valve uses thermionic emission of electrons from a heated cathode for a number of fundamental electronic functions such as signal amplification and current rectification.
The simplest vacuum tube, the diode invented in 1904 by John Ambrose Fleming, contains only a heated electron-emitting cathode and an anode. Electrons can only flow in one direction through the device—from the cathode to the anode. Adding one or more control grids within the tube enables the current between the cathode and anode to be controlled by the voltage on the grid or grids. These devices became a key component of electronic circuits for the first half of the 20th century and were crucial to the development of radio, television, radar, sound recording and reproduction, long-distance telephone networks, and analogue and early digital computers. While some applications had used earlier technologies such as the spark gap transmitter for radio or mechanical computers for computing, it was the invention of the thermionic vacuum tube that made these technologies widespread and practical, leading to the creation of electronics.
In the 1940s, the invention of semiconductor devices made it possible to produce solid-state devices, which are smaller, cheaper, and more efficient, reliable, and durable than thermionic tubes. Starting in the mid-1960s, thermionic tubes were replaced with the transistor. Thermionic tubes still have some applications for certain high-frequency amplifiers.
Computer networks and the Internet
On 11 September 1940, George Stibitz transmitted problems for his Complex Number Calculator in New York using a teletype and received the computed results back at Dartmouth College in New Hampshire. This configuration of a centralized computer (mainframe) with remote dumb terminals remained popular well into the 1970s. In the 1960s, Paul Baran and, independently, Donald Davies started to investigate packet switching, a technology that sends a message in portions to its destination asynchronously without passing it through a centralized mainframe. A four-node network emerged on 5 December 1969, constituting the beginnings of the ARPANET, which by 1981 had grown to 213 nodes. ARPANET eventually merged with other networks to form the Internet. While Internet development was a focus of the Internet Engineering Task Force (IETF) who published a series of Request for Comments documents, other networking advancements occurred in industrial laboratories, such as the local area network (LAN) developments of Ethernet (1983), Token Ring (1984) and Star network topology.
Growth of transmission capacity
The effective capacity to exchange information worldwide through two-way telecommunication networks grew from 281 petabytes (PB) of optimally compressed information in 1986 to 471 PB in 1993 to 2.2 exabytes (EB) in 2000 to 65 EB in 2007. This is the informational equivalent of two newspaper pages per person per day in 1986, and six entire newspapers per person per day by 2007. Given this growth, telecommunications play an increasingly important role in the world economy and the global telecommunications industry was about a $4.7 trillion sector in 2012. The service revenue of the global telecommunications industry was estimated to be $1.5 trillion in 2010, corresponding to 2.4% of the world's gross domestic product (GDP).
Technical concepts
Modern telecommunication is founded on a series of key concepts that experienced progressive development and refinement in a period of well over a century:
Basic elements
Telecommunication technologies may primarily be divided into wired and wireless methods. Overall, a basic telecommunication system consists of three main parts that are always present in some form or another:
A transmitter that takes information and converts it to a signal
A transmission medium, also called the physical channel, that carries the signal (e.g., the "free space channel")
A receiver that takes the signal from the channel and converts it back into usable information for the recipient
In a radio broadcasting station, the station's large power amplifier is the transmitter and the broadcasting antenna is the interface between the power amplifier and the free space channel. The free space channel is the transmission medium and the receiver's antenna is the interface between the free space channel and the receiver. Next, the radio receiver is the destination of the radio signal, where it is converted from electricity to sound.
Telecommunication systems are occasionally "duplex" (two-way systems) with a single box of electronics working as both the transmitter and a receiver, or a transceiver (e.g., a mobile phone). The transmission electronics and the receiver electronics within a transceiver are quite independent of one another. This can be explained by the fact that radio transmitters contain power amplifiers that operate with electrical powers measured in watts or kilowatts, but radio receivers deal with radio powers measured in microwatts or nanowatts. Hence, transceivers have to be carefully designed and built to isolate their high-power circuitry and their low-power circuitry from each other to avoid interference.
Telecommunication over fixed lines is called point-to-point communication because it occurs between a transmitter and a receiver. Telecommunication through radio broadcasts is called broadcast communication because it occurs between a powerful transmitter and numerous low-power but sensitive radio receivers.
Telecommunications in which multiple transmitters and multiple receivers have been designed to cooperate and share the same physical channel are called multiplex systems. The sharing of physical channels using multiplexing often results in significant cost reduction. Multiplexed systems are laid out in telecommunication networks and multiplexed signals are switched at nodes through to the correct destination terminal receiver.
Analogue versus digital communications
Communications can be encoded as analogue or digital signals, which may in turn be carried by analogue or digital communication systems. Analogue signals vary continuously with respect to the information, while digital signals encode information as a set of discrete values (e.g., a set of ones and zeroes). During propagation and reception, information contained in analogue signals is degraded by undesirable noise. Commonly, the noise in a communication system can be expressed as adding or subtracting from the desirable signal via a random process. This form of noise is called additive noise, with the understanding that the noise can be negative or positive at different instances.
Unless the additive noise disturbance exceeds a certain threshold, the information contained in digital signals will remain intact. Their resistance to noise represents a key advantage of digital signals over analogue signals. However, digital systems fail catastrophically when noise exceeds the system's ability to autocorrect. On the other hand, analogue systems fail gracefully: as noise increases, the signal becomes progressively more degraded but still usable. Also, digital transmission of continuous data unavoidably adds quantization noise to the output. This can be reduced, but not eliminated, only at the expense of increasing the channel bandwidth requirement.
Communication channels
The term channel has two different meanings. In one meaning, a channel is the physical medium that carries a signal between the transmitter and the receiver. Examples of this include the atmosphere for sound communications, glass optical fibres for some kinds of optical communications, coaxial cables for communications by way of the voltages and electric currents in them, and free space for communications using visible light, infrared waves, ultraviolet light, and radio waves. Coaxial cable types are classified by RG type or radio guide, terminology derived from World War II. The various RG designations are used to classify the specific signal transmission applications. This last channel is called the free space channel. The sending of radio waves from one place to another has nothing to do with the presence or absence of an atmosphere between the two. Radio waves travel through a perfect vacuum just as easily as they travel through air, fog, clouds, or any other kind of gas.
The other meaning of the term channel in telecommunications is seen in the phrase communications channel, which is a subdivision of a transmission medium so that it can be used to send multiple streams of information simultaneously. For example, one radio station can broadcast radio waves into free space at frequencies in the neighbourhood of 94.5 MHz (megahertz) while another radio station can simultaneously broadcast radio waves at frequencies in the neighbourhood of 96.1 MHz. Each radio station would transmit radio waves over a frequency bandwidth of about 180 kHz (kilohertz), centred at frequencies such as the above, which are called the "carrier frequencies". Each station in this example is separated from its adjacent stations by 200 kHz, and the difference between 200 kHz and 180 kHz (20 kHz) is an engineering allowance for the imperfections in the communication system.
In the example above, the free space channel has been divided into communications channels according to frequencies, and each channel is assigned a separate frequency bandwidth in which to broadcast radio waves. This system of dividing the medium into channels according to frequency is called frequency-division multiplexing. Another term for the same concept is wavelength-division multiplexing, which is more commonly used in optical communications when multiple transmitters share the same physical medium.
Another way of dividing a communications medium into channels is to allocate each sender a recurring segment of time (a time slot, for example, 20 milliseconds out of each second), and to allow each sender to send messages only within its own time slot. This method of dividing the medium into communication channels is called time-division multiplexing (TDM), and is used in optical fibre communication. Some radio communication systems use TDM within an allocated FDM channel. Hence, these systems use a hybrid of TDM and FDM.
Modulation
The shaping of a signal to convey information is known as modulation. Modulation can be used to represent a digital message as an analogue waveform. This is commonly called "keying"—a term derived from the older use of Morse Code in telecommunications—and several keying techniques exist (these include phase-shift keying, frequency-shift keying, and amplitude-shift keying). The Bluetooth system, for example, uses phase-shift keying to exchange information between various devices. In addition, there are combinations of phase-shift keying and amplitude-shift keying which is called (in the jargon of the field) quadrature amplitude modulation (QAM) that are used in high-capacity digital radio communication systems.
Modulation can also be used to transmit the information of low-frequency analogue signals at higher frequencies. This is helpful because low-frequency analogue signals cannot be effectively transmitted over free space. Hence the information from a low-frequency analogue signal must be impressed into a higher-frequency signal (known as the carrier wave) before transmission. There are several different modulation schemes available to achieve this [two of the most basic being amplitude modulation (AM) and frequency modulation (FM)]. An example of this process is a disc jockey's voice being impressed into a 96 MHz carrier wave using frequency modulation (the voice would then be received on a radio as the channel 96 FM). In addition, modulation has the advantage that it may use frequency division multiplexing (FDM).
Telecommunication networks
A telecommunications network is a collection of transmitters, receivers, and communications channels that send messages to one another. Some digital communications networks contain one or more routers that work together to transmit information to the correct user. An analogue communications network consists of one or more switches that establish a connection between two or more users. For both types of networks, repeaters may be necessary to amplify or recreate the signal when it is being transmitted over long distances. This is to combat attenuation that can render the signal indistinguishable from the noise. Another advantage of digital systems over analogue is that their output is easier to store in memory, i.e., two voltage states (high and low) are easier to store than a continuous range of states.
Societal impact
Telecommunication has a significant social, cultural and economic impact on modern society. In 2008, estimates placed the telecommunication industry's revenue at US$4.7 trillion or just under three per cent of the gross world product (official exchange rate). Several following sections discuss the impact of telecommunication on society.
Microeconomics
On the microeconomic scale, companies have used telecommunications to help build global business empires. This is self-evident in the case of online retailer Amazon.com but, according to academic Edward Lenert, even the conventional retailer Walmart has benefited from better telecommunication infrastructure compared to its competitors. In cities throughout the world, home owners use their telephones to order and arrange a variety of home services ranging from pizza deliveries to electricians. Even relatively poor communities have been noted to use telecommunication to their advantage. In Bangladesh's Narsingdi District, isolated villagers use cellular phones to speak directly to wholesalers and arrange a better price for their goods. In Côte d'Ivoire, coffee growers share mobile phones to follow hourly variations in coffee prices and sell at the best price.
Macroeconomics
On the macroeconomic scale, Lars-Hendrik Röller and Leonard Waverman suggested a causal link between good telecommunication infrastructure and economic growth. Few dispute the existence of a correlation although some argue it is wrong to view the relationship as causal.
Because of the economic benefits of good telecommunication infrastructure, there is increasing worry about the inequitable access to telecommunication services amongst various countries of the world—this is known as the digital divide. A 2003 survey by the International Telecommunication Union (ITU) revealed that roughly a third of countries have fewer than one mobile subscription for every 20 people and one-third of countries have fewer than one land-line telephone subscription for every 20 people. In terms of Internet access, roughly half of all countries have fewer than one out of 20 people with Internet access. From this information, as well as educational data, the ITU was able to compile an index that measures the overall ability of citizens to access and use information and communication technologies. Using this measure, Sweden, Denmark and Iceland received the highest ranking while the African countries Niger, Burkina Faso and Mali received the lowest.
Social impact
Telecommunication has played a significant role in social relationships. Nevertheless, devices like the telephone system were originally advertised with an emphasis on the practical dimensions of the device (such as the ability to conduct business or order home services) as opposed to the social dimensions. It was not until the late 1920s and 1930s that the social dimensions of the device became a prominent theme in telephone advertisements. New promotions started appealing to consumers' emotions, stressing the importance of social conversations and staying connected to family and friends.
Since then the role that telecommunications has played in social relations has become increasingly important. In recent years, the popularity of social networking sites has increased dramatically. These sites allow users to communicate with each other as well as post photographs, events and profiles for others to see. The profiles can list a person's age, interests, sexual preference and relationship status. In this way, these sites can play important role in everything from organising social engagements to courtship.
Prior to social networking sites, technologies like short message service (SMS) and the telephone also had a significant impact on social interactions. In 2000, market research group Ipsos MORI reported that 81% of 15- to 24-year-old SMS users in the United Kingdom had used the service to coordinate social arrangements and 42% to flirt.
Entertainment, news, and advertising
In cultural terms, telecommunication has increased the public's ability to access music and film. With television, people can watch films they have not seen before in their own home without having to travel to the video store or cinema. With radio and the Internet, people can listen to music they have not heard before without having to travel to the music store.
Telecommunication has also transformed the way people receive their news. A 2006 survey (right table) of slightly more than 3,000 Americans by the non-profit Pew Internet and American Life Project in the United States the majority specified television or radio over newspapers.
Telecommunication has had an equally significant impact on advertising. TNS Media Intelligence reported that in 2007, 58% of advertising expenditure in the United States was spent on media that depend upon telecommunication.
Regulation
Many countries have enacted legislation which conforms to the International Telecommunication Regulations established by the International Telecommunication Union (ITU), which is the "leading UN agency for information and communication technology issues". In 1947, at the Atlantic City Conference, the ITU decided to "afford international protection to all frequencies registered in a new international frequency list and used in conformity with the Radio Regulation". According to the ITU's Radio Regulations adopted in Atlantic City, all frequencies referenced in the International Frequency Registration Board, examined by the board and registered on the International Frequency List "shall have the right to international protection from harmful interference".
From a global perspective, there have been political debates and legislation regarding the management of telecommunication and broadcasting. The history of broadcasting discusses some debates in relation to balancing conventional communication such as printing and telecommunication such as radio broadcasting. The onset of World War II brought on the first explosion of international broadcasting propaganda. Countries, their governments, insurgents, terrorists, and militiamen have all used telecommunication and broadcasting techniques to promote propaganda. Patriotic propaganda for political movements and colonization started the mid-1930s. In 1936, the BBC broadcast propaganda to the Arab World to partly counter similar broadcasts from Italy, which also had colonial interests in North Africa. Modern political debates in telecommunication include the reclassification of broadband Internet service as a telecommunications service (also called net neutrality), regulation of phone spam, and expanding affordable broadband access.
Modern media
Worldwide equipment sales
According to data collected by Gartner and Ars Technica sales of main consumer's telecommunication equipment worldwide in millions of units was:
Telephone
In a telephone network, the caller is connected to the person to whom they wish to talk by switches at various telephone exchanges. The switches form an electrical connection between the two users and the setting of these switches is determined electronically when the caller dials the number. Once the connection is made, the caller's voice is transformed to an electrical signal using a small microphone in the caller's handset. This electrical signal is then sent through the network to the user at the other end where it is transformed back into sound by a small speaker in that person's handset.
, the landline telephones in most residential homes are analogue—that is, the speaker's voice directly determines the signal's voltage. Although short-distance calls may be handled from end-to-end as analogue signals, increasingly telephone service providers are transparently converting the signals to digital signals for transmission. The advantage of this is that digitized voice data can travel side by side with data from the Internet and can be perfectly reproduced in long-distance communication (as opposed to analogue signals that are inevitably impacted by noise).
Mobile phones have had a significant impact on telephone networks. Mobile phone subscriptions now outnumber fixed-line subscriptions in many markets. Sales of mobile phones in 2005 totalled 816.6 million with that figure being almost equally shared amongst the markets of Asia/Pacific (204 m), Western Europe (164 m), CEMEA (Central Europe, the Middle East and Africa) (153.5 m), North America (148 m) and Latin America (102 m). In terms of new subscriptions over the five years from 1999, Africa has outpaced other markets with 58.2% growth. Increasingly these phones are being serviced by systems where the voice content is transmitted digitally such as GSM or W-CDMA with many markets choosing to deprecate analog systems such as AMPS.
There have also been dramatic changes in telephone communication behind the scenes. Starting with the operation of TAT-8 in 1988, the 1990s saw the widespread adoption of systems based on optical fibres. The benefit of communicating with optical fibres is that they offer a drastic increase in data capacity. TAT-8 itself was able to carry 10 times as many telephone calls as the last copper cable laid at that time and today's optical fibre cables are able to carry 25 times as many telephone calls as TAT-8. This increase in data capacity is due to several factors: First, optical fibres are physically much smaller than competing technologies. Second, they do not suffer from crosstalk which means several hundred of them can be easily bundled together in a single cable. Lastly, improvements in multiplexing have led to an exponential growth in the data capacity of a single fibre.
Assisting communication across many modern optical fibre networks is a protocol known as Asynchronous Transfer Mode (ATM). The ATM protocol allows for the side-by-side data transmission mentioned in the second paragraph. It is suitable for public telephone networks because it establishes a pathway for data through the network and associates a traffic contract with that pathway. The traffic contract is essentially an agreement between the client and the network about how the network is to handle the data; if the network cannot meet the conditions of the traffic contract it does not accept the connection. This is important because telephone calls can negotiate a contract so as to guarantee themselves a constant bit rate, something that will ensure a caller's voice is not delayed in parts or cut off completely. There are competitors to ATM, such as Multiprotocol Label Switching (MPLS), that perform a similar task and are expected to supplant ATM in the future.
Radio and television
In a broadcast system, the central high-powered broadcast tower transmits a high-frequency electromagnetic wave to numerous low-powered receivers. The high-frequency wave sent by the tower is modulated with a signal containing visual or audio information. The receiver is then tuned so as to pick up the high-frequency wave and a demodulator is used to retrieve the signal containing the visual or audio information. The broadcast signal can be either analogue (signal is varied continuously with respect to the information) or digital (information is encoded as a set of discrete values).
The broadcast media industry is at a critical turning point in its development, with many countries moving from analogue to digital broadcasts. This move is made possible by the production of cheaper, faster and more capable integrated circuits. The chief advantage of digital broadcasts is that they prevent a number of complaints common to traditional analogue broadcasts. For television, this includes the elimination of problems such as snowy pictures, ghosting and other distortion. These occur because of the nature of analogue transmission, which means that perturbations due to noise will be evident in the final output. Digital transmission overcomes this problem because digital signals are reduced to discrete values upon reception and hence small perturbations do not affect the final output. In a simplified example, if a binary message 1011 was transmitted with signal amplitudes [1.0 0.0 1.0 1.0] and received with signal amplitudes [0.9 0.2 1.1 0.9] it would still decode to the binary message 1011— a perfect reproduction of what was sent. From this example, a problem with digital transmissions can also be seen in that if the noise is great enough it can significantly alter the decoded message. Using forward error correction a receiver can correct a handful of bit errors in the resulting message but too much noise will lead to incomprehensible output and hence a breakdown of the transmission.
In digital television broadcasting, there are three competing standards that are likely to be adopted worldwide. These are the ATSC, DVB and ISDB standards; the adoption of these standards thus far is presented in the captioned map. All three standards use MPEG-2 for video compression. ATSC uses Dolby Digital AC-3 for audio compression, ISDB uses Advanced Audio Coding (MPEG-2 Part 7) and DVB has no standard for audio compression but typically uses MPEG-1 Part 3 Layer 2. The choice of modulation also varies between the schemes. In digital audio broadcasting, standards are much more unified with practically all countries choosing to adopt the Digital Audio Broadcasting standard (also known as the Eureka 147 standard). The exception is the United States which has chosen to adopt HD Radio. HD Radio, unlike Eureka 147, is based upon a transmission method known as in-band on-channel transmission that allows digital information to piggyback on normal AM or FM analog transmissions.
However, despite the pending switch to digital, analog television remains being transmitted in most countries. An exception is the United States that ended analog television transmission (by all but the very low-power TV stations) on 12 June 2009 after twice delaying the switchover deadline. Kenya also ended analog television transmission in December 2014 after multiple delays. For analogue television, there were three standards in use for broadcasting colour TV (see a map on adoption here). These are known as PAL (German designed), NTSC (American designed), and SECAM (French-designed). For analogue radio, the switch to digital radio is made more difficult by the higher cost of digital receivers. The choice of modulation for analogue radio is typically between amplitude (AM) or frequency modulation (FM). To achieve stereo playback, an amplitude modulated subcarrier is used for stereo FM, and quadrature amplitude modulation is used for stereo AM or C-QUAM.
Internet
The Internet is a worldwide network of computers and computer networks that communicate with each other using the Internet Protocol (IP). Any computer on the Internet has a unique IP address that can be used by other computers to route information to it. Hence, any computer on the Internet can send a message to any other computer using its IP address. These messages carry with them the originating computer's IP address allowing for two-way communication. The Internet is thus an exchange of messages between computers.
It is estimated that 51% of the information flowing through two-way telecommunications networks in the year 2000 were flowing through the Internet (most of the rest (42%) through the landline telephone). By 2007 the Internet clearly dominated and captured 97% of all the information in telecommunication networks (most of the rest (2%) through mobile phones). , an estimated 21.9% of the world population has access to the Internet with the highest access rates (measured as a percentage of the population) in North America (73.6%), Oceania/Australia (59.5%) and Europe (48.1%). In terms of broadband access, Iceland (26.7%), South Korea (25.4%) and the Netherlands (25.3%) led the world.
The Internet works in part because of protocols that govern how the computers and routers communicate with each other. The nature of computer network communication lends itself to a layered approach where individual protocols in the protocol stack run more-or-less independently of other protocols. This allows lower-level protocols to be customized for the network situation while not changing the way higher-level protocols operate. A practical example of why this is important is because it allows a web browser to run the same code regardless of whether the computer it is running on is connected to the Internet through an Ethernet or Wi-Fi connection. Protocols are often talked about in terms of their place in the OSI reference model (pictured on the right), which emerged in 1983 as the first step in an unsuccessful attempt to build a universally adopted networking protocol suite.
For the Internet, the physical medium and data link protocol can vary several times as packets traverse the globe. This is because the Internet places no constraints on what physical medium or data link protocol is used. This leads to the adoption of media and protocols that best suit the local network situation. In practice, most intercontinental communication will use the Asynchronous Transfer Mode (ATM) protocol (or a modern equivalent) on top of optic fibre. This is because for most intercontinental communication the Internet shares the same infrastructure as the public switched telephone network.
At the network layer, things become standardized with the Internet Protocol (IP) being adopted for logical addressing. For the World Wide Web, these IP addresses are derived from the human-readable form using the Domain Name System (e.g., 72.14.207.99 is derived from Google.com). At the moment, the most widely used version of the Internet Protocol is version four but a move to version six is imminent.
At the transport layer, most communication adopts either the Transmission Control Protocol (TCP) or the User Datagram Protocol (UDP). TCP is used when it is essential every message sent is received by the other computer whereas UDP is used when it is merely desirable. With TCP, packets are retransmitted if they are lost and placed in order before they are presented to higher layers. With UDP, packets are not ordered nor retransmitted if lost. Both TCP and UDP packets carry port numbers with them to specify what application or process the packet should be handled by. Because certain application-level protocols use certain ports, network administrators can manipulate traffic to suit particular requirements. Examples are to restrict Internet access by blocking the traffic destined for a particular port or to affect the performance of certain applications by assigning priority.
Above the transport layer, there are certain protocols that are sometimes used and loosely fit in the session and presentation layers, most notably the Secure Sockets Layer (SSL) and Transport Layer Security (TLS) protocols. These protocols ensure that data transferred between two parties remains completely confidential. Finally, at the application layer, are many of the protocols Internet users would be familiar with such as HTTP (web browsing), POP3 (e-mail), FTP (file transfer), IRC (Internet chat), BitTorrent (file sharing) and XMPP (instant messaging).
Voice over Internet Protocol (VoIP) allows data packets to be used for synchronous voice communications. The data packets are marked as voice-type packets and can be prioritized by the network administrators so that the real-time, synchronous conversation is less subject to contention with other types of data traffic which can be delayed (i.e., file transfer or email) or buffered in advance (i.e., audio and video) without detriment. That prioritization is fine when the network has sufficient capacity for all the VoIP calls taking place at the same time and the network is enabled for prioritization, i.e., a private corporate-style network, but the Internet is not generally managed in this way and so there can be a big difference in the quality of VoIP calls over a private network and over the public Internet.
Local area networks and wide area networks
Despite the growth of the Internet, the characteristics of local area networks (LANs)—computer networks that do not extend beyond a few kilometres—remain distinct. This is because networks on this scale do not require all the features associated with larger networks and are often more cost-effective and efficient without them. When they are not connected with the Internet, they also have the advantages of privacy and security. However, purposefully lacking a direct connection to the Internet does not provide assured protection from hackers, military forces, or economic powers. These threats exist if there are any methods for connecting remotely to the LAN.
Wide area networks (WANs) are private computer networks that may extend for thousands of kilometres. Once again, some of their advantages include privacy and security. Prime users of private LANs and WANs include armed forces and intelligence agencies that must keep their information secure and secret.
In the mid-1980s, several sets of communication protocols emerged to fill the gaps between the data-link layer and the application layer of the OSI reference model. These included AppleTalk, IPX, and NetBIOS with the dominant protocol set during the early 1990s being IPX due to its popularity with MS-DOS users. TCP/IP existed at this point, but it was typically only used by large government and research facilities.
As the Internet grew in popularity and its traffic was required to be routed into private networks, the TCP/IP protocols replaced existing local area network technologies. Additional technologies, such as DHCP, allowed TCP/IP-based computers to self-configure in the network. Such functions also existed in the AppleTalk/ IPX/ NetBIOS protocol sets.
Whereas Asynchronous Transfer Mode (ATM) or Multiprotocol Label Switching (MPLS) are typical data-link protocols for larger networks such as WANs; Ethernet and Token Ring are typical data-link protocols for LANs. These protocols differ from the former protocols in that they are simpler, e.g., they omit features such as quality of service guarantees, and offer medium access control. Both of these differences allow for more economical systems.
Despite the modest popularity of Token Ring in the 1980s and 1990s, virtually all LANs now use either wired or wireless Ethernet facilities. At the physical layer, most wired Ethernet implementations use copper twisted-pair cables (including the common 10BASE-T networks). However, some early implementations used heavier coaxial cables and some recent implementations (especially high-speed ones) use optical fibres. When optic fibres are used, the distinction must be made between multimode fibres and single-mode fibres. Multimode fibres can be thought of as thicker optical fibres that are cheaper to manufacture devices for, but that suffer from less usable bandwidth and worse attenuation—implying poorer long-distance performance.
| Technology | Media and communication | null |
930437 | https://en.wikipedia.org/wiki/Sungrazing%20comet | Sungrazing comet | A sungrazing comet is a comet that passes extremely close to the Sun at perihelion – sometimes within a few thousand kilometres of the Sun's surface. Although small sungrazers can completely evaporate during such a close approach to the Sun, larger sungrazers can survive many perihelion passages. However, the strong evaporation and tidal forces they experience often lead to their fragmentation.
Up until the 1880s, it was thought that all bright comets near the Sun were the repeated return of a single sungrazing comet. Then, German astronomer Heinrich Kreutz and American astronomer Daniel Kirkwood determined that, instead of the return of the same comet, each appearance was a different comet, but each were related to a group of comets that had separated from each other at an earlier passage near the Sun (at perihelion). Very little was known about the population of sungrazing comets until 1979 when coronagraphic observations allowed the detection of sungrazers. As of October 21, 2017, there are 1495 known comets that come within ~12 solar radii (~0.055 AU). This accounts for nearly one third of all comets. Most of these objects vaporize during their close approach, but a comet with a nucleus radius larger than 2–3 km is likely to survive the perihelion passage with a final radius of ~1 km.
Sungrazer comets were some of the earliest observed comets because they can appear very bright. Some are even considered Great Comets. The close passage of a comet to the Sun will brighten the comet not only because of the reflection off the comet nucleus when it is closer to the Sun, but the Sun also vaporizes a large amount of gas from the comet and the gas reflects more light. This extreme brightening will allow for possible naked eye observations from Earth depending on how volatile the gases are and if the comet is large enough to survive perihelion. These comets provide a useful tool for understanding the composition of comets as we observe the outgassing activity and they also offer a way to probe the effects solar radiation has on other Solar System bodies.
History of sungrazers
Pre-19th century
One of the first comets to have its orbit computed was the sungrazing comet (and Great Comet) of 1680, now designated C/1680 V1. It was observed by Isaac Newton and he published the orbit results in 1687. Later, in 1699, Jacques Cassini proposed that comets could have relatively short orbital periods and that C/1680 V1 was the same as a comet observed by Tycho Brahe in 1577, but in 1705 Edmond Halley determined that the difference between the perihelion distances of the two comets was too great for them to be the same object. However, this marked the first time that it was hypothesized that Great Comets were related or perhaps the same comet. Later, Johann Franz Encke computed the orbit of C/1680 V1 and found a period of approximately 9000 years, leading him to conclude that Cassini's theory of short period sungrazers was flawed. C/1680 V1 had the smallest measured perihelion distance until the observation in 1826 of comet C/1826 U1.
19th century
Advances were made in understanding sungrazing comets in the 19th century with the Great Comets of 1843, C/1880 C1, and 1882. C/1880 C1 and C/1843 D1 had very similar appearances and also resembled the Great Comet of 1106, therefore Daniel Kirkwood proposed that C/1880 C1 and C/1843 D1 were separate fragments of the same object. He also hypothesized that the parent body was a comet seen by Aristotle and Ephorus in 371 BC because there was a supposed claim that Ephorus witnessed the comet splitting after perihelion.
Comet C/1882 R1 appeared only two years after the previously observed sungrazer so this convinced astronomers that these bright comets were not all the same object. Some astronomers theorized that the comet might pass through a resisting medium near the Sun and that would shorten its period. When astronomers observed C/1882 R1, they measured the period before and after perihelion and saw no shortening in the period which disproved the theory. After perihelion this object was also seen to split into several fragments and therefore Kirkwood's theory of these comets coming from a parent body seemed like a good explanation.
In an attempt to link the 1843 and 1880 comets to the comet in 1106 and 371 BC, Kreutz measured the fragments of the 1882 comet and determined that it was likely a fragment of the 1106 comet. He then designated that all sungrazing comets with similar orbital characteristics as these few comets would be part of the Kreutz Group.
The 19th century also provided the first spectrum taken of a comet near the Sun which was taken by Finlay & Elkin in 1882. Later the spectrum was analyzed and Fe and Ni spectral lines were confirmed.
20th century
The first sungrazing comet observed in the 20th century was in 1945 and then between 1960 and 1970 five sungrazing comets were seen (C/1961 O1, C/1962 C1, C/1963 R1, C/1965 S1, and C/1970 K1). The 1965 comet (Comet Ikeya-Seki) allowed for measurements of spectral emission lines and several elements were detected including Iron, marking this the first comet since the Great Comet of 1882 to show this feature. Other emission lines included K, Ca, Ca+, Cr, Co, Mn, Ni, Cu, and V. Comet Ikeya-Seki also led to separating the Kreutz sungrazers into two subgroups by Brian Marsden in 1967. One subgroup appears to have the 1106 comet as the parent body and members are fragments of that comet, while the other group have similar dynamics but no confirmed parent body associated with it.
Coronagraphic observations
The 20th century greatly impacted sungrazing comet research with the launch of coronagraphic telescopes including Solwind, SMM, and SOHO. Until this point, sungrazing comets were only seen with the naked eye but with the coronagraphic telescopes many sungrazers were observed that were much smaller and very few have survived perihelion passage. The comets observed by Solwind and SMM from 1981 to 1989 had visual magnitudes from about -2.5 to +6 which is much fainter than Comet Ikeya-Seki with a visual magnitude of about -10.
In 1987 and 1988 it was first observed by SMM that there could be pairs of sungrazing comets that can appear within very short time periods ranging from a half of a day up to about two weeks. Calculations were made to determine that the pairs were part of the same parent body but broke apart at tens of AU from the Sun. The breakup velocities were only on the order of a few meters per second which is comparable to the speed of rotation for these comets. This led to the conclusion that these comets break from tidal forces and that comets C/1882 R1, C/1965 S1, and C/1963 R1 probably broke off from the Great Comet of 1106.
Coronagraphs allowed for measuring the properties of the comet as it reached very close to the Sun. It was noted that sungrazing comets tend to peak in brightness at a distance of about 12.3 solar radii or 11.2 solar radii. It is thought that this variation stems from a difference in dust composition. Another small peak in brightness has been found at about 7 solar radii from the sun and it is possibly due to a fragmentation of the comet nucleus. An alternative explanation is that the brightness peak at 12 solar radii comes from the sublimation of amorphous olivines and the peak at 11.2 solar radii is from the sublimation of crystalline olivines. The peak at 7 solar radii could then be the sublimation of pyroxene.
Sungrazing groups
Kreutz Sungrazers
The most famous sungrazers are the Kreutz Sungrazers, which all originate from one giant comet that broke up into many smaller comets during its first passage through the inner Solar System. An extremely bright comet seen by Aristotle and Ephorus in 371 BC is a possible candidate for this parent comet.
The Great Comets of 1843 and 1882, Comet Ikeya–Seki in 1965 and C/2011 W3 (Lovejoy) in 2011 were all fragments of the original comet. Each of these four was briefly bright enough to be visible in the daytime sky, next to the Sun, 1882's comet outshining even the full moon.
In 1979, C/1979 Q1 (Solwind) was the first sungrazer to be spotted by US satellite P78-1, in coronagraphs taken on 30 and 31 Aug 1979.
Apart from Comet Lovejoy, none of the sungrazers seen by SOHO has survived its perihelion passage; some may have plunged into the Sun itself, but most are likely to have simply evaporated away completely.
Other sungrazers
About 83% of the sungrazers observed with SOHO are members of the Kreutz group. The other 17% contains some sporadic sungrazers, but three other related groups of comets have been identified among them: the Kracht, Marsden and Meyer groups. The Marsden and Kracht groups both appear to be related to Comet 96P/Machholz. These comets have also been linked to several meteor streams, including the Daytime Arietids, the delta Aquariids, and the Quadrantids. Linked comet orbits suggest that both Marsden and Kracht groups have a small period, on the order of five years, but the Meyer group may have intermediate- or long-period orbits. The Meyer group comets are typically small, faint, and never have tails. The Great Comet of 1680 was a sungrazer and while used by Newton to verify Kepler's equations on orbital motion, it was not a member of any larger groups. However, comet C/2012 S1 (ISON), which disintegrated shortly before perihelion, had orbital elements similar to the Great Comet of 1680 and could be a second member of the group.
Origin of sungrazing comets
Studies show that for comets with high orbital inclinations and perihelion distances of less than about 2 astronomical units, the cumulative effect of gravitational perturbations over many orbits is adequate to reduce the perihelion distance to very small values. One study has suggested that Comet Hale–Bopp has about a 15% chance of eventually becoming a sungrazer.
Role in solar astronomy
The motion of tails of sungrazers that survive perihelion (such as Comet Lovejoy) can provide solar astronomers with information about the structure of the solar corona, particularly the detailed magnetic structure.
| Physical sciences | Planetary science | Astronomy |
931370 | https://en.wikipedia.org/wiki/Tropical%20rainforest | Tropical rainforest | Tropical rainforests are dense and warm rainforests with high rainfall typically found between 10° north and south of the Equator. They are a subset of the tropical forest biome that occurs roughly within the 28° latitudes (in the torrid zone between the Tropic of Cancer and Tropic of Capricorn). Tropical rainforests are a type of tropical moist broadleaf forest, that includes the more extensive seasonal tropical forests. True rainforests usually occur in tropical rainforest climates where no dry season occurs; all months have an average precipitation of at least . Seasonal tropical forests with tropical monsoon or savanna climates are sometimes included in the broader definition.
Tropical rainforests ecosystems are distinguished by their consistent, high temperatures, exceeding monthly, and substantial annual rainfall. The abundant rainfall results in nutrient-poor, leached soils, which profoundly affect the flora and fauna adapted to these conditions. These rainforests are renowned for their significant biodiversity. They are home to 40–75% of all species globally, including half of the world's animal and plant species, and two-thirds of all flowering plant species. Their dense insect population and variety of trees and higher plants are notable. Described as the "world's largest pharmacy", over a quarter of natural medicines have been discovered in them. However, tropical rainforests are threatened by human activities, such as logging and agricultural expansion, leading to habitat fragmentation and loss.
The structure of a tropical rainforest is stratified into layers, each hosting unique ecosystems. These include the emergent layer with towering trees, the densely populated canopy layer, the understory layer rich in wildlife, and the forest floor, which is sparse due to low light penetration. The soil is characteristically nutrient-poor and acidic. Tropical rainforests have a long history of ecological succession, influenced by natural events and human activities. They are crucial for global ecological functions, including carbon sequestration and climate regulation. Many indigenous peoples around the world have inhabited rainforests for millennia, relying on them for sustenance and shelter, but face challenges from modern economic activities.
Conservation efforts are diverse, focusing on both preservation and sustainable management. International policies, such as the Reducing Emissions from Deforestation and Forest Degradation (REDD and REDD+) programs, aim to curb deforestation and forest degradation. Despite these efforts, tropical rainforests continue to face significant threats from deforestation and climate change, highlighting the ongoing challenge of balancing conservation with human development needs.
Overview
Tropical rainforests are hot and wet. Mean monthly temperatures exceed during all months of the year. Average annual rainfall is no less than and can exceed although it typically lies between and . This high level of precipitation often results in poor soils due to leaching of soluble nutrients in the ground.
Tropical rainforests exhibit high levels of biodiversity. Around 40% to 75% of all biotic species are indigenous to the rainforests. Rainforests are home to half of all the living animal and plant species on the planet. Two-thirds of all flowering plants can be found in rainforests. A single hectare of rainforest may contain 42,000 different species of insect, up to 807 trees of 313 species and 1,500 species of higher plants. Tropical rainforests have been called the "world's largest pharmacy", because over one quarter of natural medicines have been discovered within them. It is likely that there may be many millions of species of plants, insects and microorganisms still undiscovered in tropical rainforests.
Tropical rainforests are among the most threatened ecosystems globally due to large-scale fragmentation as a result of human activity. Habitat fragmentation caused by geological processes such as volcanism and climate change occurred in the past, and have been identified as important drivers of speciation. However, fast human driven habitat destruction is suspected to be one of the major causes of species extinction. Tropical rain forests have been subjected to heavy logging and agricultural clearance throughout the 20th century, and the area covered by rainforests around the world is rapidly shrinking.
History
Tropical rainforests have existed on earth for hundreds of millions of years. Most tropical rainforests today are on fragments of the Mesozoic era supercontinent of Gondwana. The separation of the landmass resulted in a great loss of amphibian diversity while at the same time the drier climate spurred the diversification of reptiles. The division left tropical rainforests located in five major regions of the world: tropical America, Africa, Southeast Asia, Madagascar, and New Guinea, with smaller outliers in Australia. However, the specifics of the origin of rainforests remain uncertain due to an incomplete fossil record.
Other types of tropical forest
Several biomes may appear similar-to, or merge via ecotones with, tropical rainforest:
Moist seasonal tropical forest
Moist seasonal tropical forests receive high overall rainfall with a warm summer wet season and a cooler winter dry season. These forests usually fall under tropical monsoon or tropical savanna climates. Some trees in these forests drop some or all of their leaves during the winter dry season, thus they are sometimes called "tropical mixed forest". They are found in parts of South America, in Central America and around the Caribbean, in coastal West Africa, parts of the Indian subcontinent, and across much of Indochina.
Montane rainforests
These are found in cooler-climate mountainous areas, becoming known as cloud forests at higher elevations. Depending on latitude, the lower limit of montane rainforests on large mountains is generally between 1500 and 2500m while the upper limit is usually from 2400 to 3300m.
Flooded rainforests
Tropical freshwater swamp forests, or "flooded forests", are found in Amazon basin (the Várzea) and elsewhere.
Forest structure
Rainforests are divided into different strata, or layers, with vegetation organized into a vertical pattern from the top of the soil to the canopy. Each layer is a unique biotic community containing different plants and animals adapted for life in that particular strata. Only the emergent layer is unique to tropical rainforests, while the others are also found in temperate rainforests.
Forest floor
The forest floor, the bottom-most layer, receives only 2% of the sunlight. Only plants adapted to low light can grow in this region. Away from riverbanks, swamps and clearings, where dense undergrowth is found, the forest floor is relatively clear of vegetation because of the low sunlight penetration. This more open quality permits the easy movement of larger animals such as: ungulates like the okapi (Okapia johnstoni), tapir (Tapirus sp.), Sumatran rhinoceros (Dicerorhinus sumatrensis), and apes like the western lowland gorilla (Gorilla gorilla), as well as many species of reptiles, amphibians, and insects. The forest floor also contains decaying plant and animal matter, which disappears quickly, because the warm, humid conditions promote rapid decay. Many forms of fungi growing here help decay the animal and plant waste.
Understory layer
The understory layer lies between the canopy and the forest floor. The understory is home to a number of birds, small mammals, insects, reptiles, and predators. Examples include leopard (Panthera pardus), poison dart frogs (Dendrobates sp.), ring-tailed coati (Nasua nasua), boa constrictor (Boa constrictor), and many species of Coleoptera. The vegetation at this layer generally consists of shade-tolerant shrubs, herbs, small trees, and large woody vines which climb into the trees to capture sunlight. Only about 5% of sunlight breaches the canopy to arrive at the understory causing true understory plants to seldom grow to 3 m (10feet). As an adaptation to these low light levels, understory plants have often evolved much larger leaves. Many seedlings that will grow to the canopy level are in the understory.
Canopy layer
The canopy is the primary layer of the forest, forming a roof over the two remaining layers. It contains the majority of the largest trees, typically 30–45 m in height. Tall, broad-leaved evergreen trees are the dominant plants. The densest areas of biodiversity are found in the forest canopy, as it often supports a rich flora of epiphytes, including orchids, bromeliads, mosses and lichens. These epiphytic plants attach to trunks and branches and obtain water and minerals from rain and debris that collects on the supporting plants. The fauna is similar to that found in the emergent layer, but more diverse. It is suggested that the total arthropod species richness of the tropical canopy might be as high as 20 million. Other species inhabiting this layer include many avian species such as the yellow-casqued wattled hornbill (Ceratogymna elata), collared sunbird (Anthreptes collaris), grey parrot (Psitacus erithacus), keel-billed toucan (Ramphastos sulfuratus), scarlet macaw (Ara macao) as well as other animals like the spider monkey (Ateles sp.), African giant swallowtail (Papilio antimachus), three-toed sloth (Bradypus tridactylus), kinkajou (Potos flavus), and tamandua (Tamandua tetradactyla).
Emergent layer
The emergent layer contains a small number of very large trees, called emergents, which grow above the general canopy, reaching heights of 45–55m, although on occasion a few species will grow to 70–80m tall. Some examples of emergents include: Hydrochorea elegans, Dipteryx panamensis, Hieronyma alchorneoides, Hymenolobium mesoamericanum, Lecythis ampla and Terminalia oblonga. These trees need to be able to withstand the hot temperatures and strong winds that occur above the canopy in some areas. Several unique faunal species inhabit this layer such as the crowned eagle (Stephanoaetus coronatus), the king colobus (Colobus polykomos), and the large flying fox (Pteropus vampyrus).
However, stratification is not always clear. Rainforests are dynamic and many changes affect the structure of the forest. Emergent or canopy trees collapse, for example, causing gaps to form. Openings in the forest canopy are widely recognized as important for the establishment and growth of rainforest trees. It is estimated that perhaps 75% of the tree species at La Selva Biological Station, Costa Rica are dependent on canopy opening for seed germination or for growth beyond sapling size, for example.
Ecology
Climates
Tropical rainforests are located around and near the equator, therefore having what is called an equatorial climate characterized by three major climatic parameters: temperature, rainfall, and dry season intensity. Other parameters that affect tropical rainforests are carbon dioxide concentrations, solar radiation, and nitrogen availability. In general, climatic patterns consist of warm temperatures and high annual rainfall. However, the abundance of rainfall changes throughout the year creating distinct moist and dry seasons. Tropical forests are classified by the amount of rainfall received each year, which has allowed ecologists to define differences in these forests that look so similar in structure. According to Holdridge's classification of tropical ecosystems, true tropical rainforests have an annual rainfall greater than 2m and annual temperature greater than 24 degrees Celsius, with a potential evapotranspiration ratio (PET) value of <0.25. However, most lowland tropical forests can be classified as tropical moist or wet forests, which differ in regards to rainfall. Tropical forest ecology- dynamics, composition, and function- are sensitive to changes in climate especially changes in rainfall.
Soils
Soil types
Soil types are highly variable in the tropics and are the result of a combination of several variables such as climate, vegetation, topographic position, parent material, and soil age. Most tropical soils are characterized by significant leaching and poor nutrients, however there are some areas that contain fertile soils. Soils throughout the tropical rainforests fall into two classifications which include the ultisols and oxisols. Ultisols are known as well weathered, acidic red clay soils, deficient in major nutrients such as calcium and potassium. Similarly, oxisols are acidic, old, typically reddish, highly weathered and leached, however are well drained compared to ultisols. The clay content of ultisols is high, making it difficult for water to penetrate and flow through. The reddish color of both soils is the result of heavy heat and moisture forming oxides of iron and aluminium, which are insoluble in water and not taken up readily by plants.
Soil chemical and physical characteristics are strongly related to above ground productivity and forest structure and dynamics. The physical properties of soil control the tree turnover rates whereas chemical properties such as available nitrogen and phosphorus control forest growth rates. The soils of the eastern and central Amazon as well as the Southeast Asian Rainforest are old and mineral poor whereas the soils of the western Amazon (Ecuador and Peru) and volcanic areas of Costa Rica are young and mineral rich. Primary productivity or wood production is highest in western Amazon and lowest in eastern Amazon which contains heavily weathered soils classified as oxisols. Additionally, Amazonian soils are greatly weathered, making them devoid of minerals like phosphorus, potassium, calcium, and magnesium, which come from rock sources. However, not all tropical rainforests occur on nutrient poor soils, but on nutrient rich floodplains and volcanic soils located in the Andean foothills, and volcanic areas of Southeast Asia, Africa, and Central America.
Oxisols, infertile, deeply weathered and severely leached, have developed on the ancient Gondwanan shields. Rapid bacterial decay prevents the accumulation of humus. The concentration of iron and aluminium oxides by the laterization process gives the oxisols a bright red color and sometimes produces minable deposits (e.g., bauxite). On younger substrates, especially of volcanic origin, tropical soils may be quite fertile.
Nutrient recycling
This high rate of decomposition is the result of phosphorus levels in the soils, precipitation, high temperatures and the extensive microorganism communities. In addition to the bacteria and other microorganisms, there are an abundance of other decomposers such as fungi and termites that aid in the process as well. Nutrient recycling is important because below ground resource availability controls the above ground biomass and community structure of tropical rainforests. These soils are typically phosphorus limited, which inhibits net primary productivity or the uptake of carbon. The soil contains microbial organisms such as bacteria, which break down leaf litter and other organic matter into inorganic forms of carbon usable by plants through a process called decomposition. During the decomposition process the microbial community is respiring, taking up oxygen and releasing carbon dioxide. The decomposition rate can be evaluated by measuring the uptake of oxygen. High temperatures and precipitation increase decomposition rate, which allows plant litter to rapidly decay in tropical regions, releasing nutrients that are immediately taken up by plants through surface or ground waters. The seasonal patterns in respiration are controlled by leaf litter fall and precipitation, the driving force moving the decomposable carbon from the litter to the soil. Respiration rates are highest early in the wet season because the recent dry season results in a large percentage of leaf litter and thus a higher percentage of organic matter being leached into the soil.
Buttress roots
A common feature of many tropical rainforests is the distinct buttress roots of trees. Instead of penetrating to deeper soil layers, buttress roots create a widespread root network at the surface for more efficient uptake of nutrients in a very nutrient poor and competitive environment. Most of the nutrients within the soil of a tropical rainforest occur near the surface because of the rapid turnover time and decomposition of organisms and leaves. Because of this, the buttress roots occur at the surface so the trees can maximize uptake and actively compete with the rapid uptake of other trees. These roots also aid in water uptake and storage, increase surface area for gas exchange, and collect leaf litter for added nutrition. Additionally, these roots reduce soil erosion and maximize nutrient acquisition during heavy rains by diverting nutrient rich water flowing down the trunk into several smaller flows while also acting as a barrier to ground flow. Also, the large surface areas these roots create provide support and stability to rainforests trees, which commonly grow to significant heights. This added stability allows these trees to withstand the impacts of severe storms, thus reducing the occurrence of fallen trees.
Forest succession
Succession is an ecological process that changes the biotic community structure over time towards a more stable, diverse community structure after an initial disturbance to the community. The initial disturbance is often a natural phenomenon or human caused event. Natural disturbances include hurricanes, volcanic eruptions, river movements or an event as small as a fallen tree that creates gaps in the forest. In tropical rainforests, these same natural disturbances have been well documented in the fossil record, and are credited with encouraging speciation and endemism. Human land use practices have led to large-scale deforestation. In many tropical countries such as Costa Rica these deforested lands have been abandoned and forests have been allowed to regenerate through ecological succession. These regenerating young successional forests are called secondary forests or second-growth forests.
Biodiversity and speciation
Tropical rainforests exhibit a vast diversity in plant and animal species. The root for this remarkable speciation has been a query of scientists and ecologists for years. A number of theories have been developed for why and how the tropics can be so diverse.
Interspecific competition
Interspecific competition results from a high density of species with similar niches in the tropics and limited resources available. Species which "lose" the competition may either become extinct or find a new niche. Direct competition will often lead to one species dominating another by some advantage, ultimately driving it to extinction. Niche partitioning is the other option for a species. This is the separation and rationing of necessary resources by utilizing different habitats, food sources, cover or general behavioral differences. A species with similar food items but different feeding times is an example of niche partitioning.
Pleistocene refugia
The theory of Pleistocene refugia was developed by Jürgen Haffer in 1969 with his article Speciation of Amazonian Forest Birds. Haffer proposed the explanation for speciation was the product of rainforest patches being separated by stretches of non-forest vegetation during the last glacial period. He called these patches of rainforest areas refuges and within these patches allopatric speciation occurred. With the end of the glacial period and increase in atmospheric humidity, rainforest began to expand and the refuges reconnected. This theory has been the subject of debate. Scientists are still skeptical of whether or not this theory is legitimate. Genetic evidence suggests speciation had occurred in certain taxa 1–2 million years ago, preceding the Pleistocene.
Human dimensions
Habitation
Tropical rainforests have harboured human life for many millennia, with many Indigenous people in South and Central America, who belong to the Indigenous peoples of the Americas, the Congo Pygmies in Central Africa, and several tribes in Southeast Asia, like the Dayak people and the Penan people in Borneo. Food resources within the forest are extremely dispersed due to the high biological diversity and what food does exist is largely restricted to the canopy and requires considerable energy to obtain. Some groups of hunter-gatherers have exploited rainforest on a seasonal basis but dwelt primarily in adjacent savanna and open forest environments where food is much more abundant. Other people described as rainforest dwellers are hunter-gatherers who subsist in large part by trading high value forest products such as hides, feathers, and honey with agricultural people living outside the forest.
Indigenous peoples
Many indigenous peoples around the world live within rainforests as hunter-gatherers, or subsist as part-time small scale farmers supplemented in large part by trading high-value forest products such as hides, feathers, and honey with agricultural people living outside the forests. Peoples have inhabited the rainforests for tens of thousands of years and have remained so elusive that only recently have some tribes been discovered. These indigenous peoples are greatly threatened by loggers in search for old-growth tropical hardwoods like Ipe, Cumaru and Wenge, and by farmers who are looking to expand their land, for cattle(meat), and soybeans, which are used to feed cattle in Europe and China. On 18 January 2007, FUNAI reported also that it had confirmed the presence of 67 different uncontacted tribes in Brazil, up from 40 in 2005. With this addition, Brazil has now overtaken the island of New Guinea as the country having the largest number of uncontacted tribes. The province of Irian Jaya or West Papua in the island of New Guinea is home to an estimated 44 uncontacted tribal groups.
The pygmy peoples are hunter-gatherer groups living in equatorial rainforests characterized by their short height (below one and a half meters, or 59inches, on average). Amongst this group are the Efe, Aka, Twa, Baka, and Mbuti people of Central Africa. However, the term pygmy is considered pejorative so many tribes prefer not to be labeled as such.
Some notable indigenous peoples of the Americas, or Amerindians, include the Huaorani, Ya̧nomamö, and Kayapo people of the Amazon. The traditional agricultural system practiced by tribes in the Amazon is based on swidden cultivation (also known as slash-and-burn or shifting cultivation) and is considered a relatively benign disturbance. In fact, when looking at the level of individual swidden plots a number of traditional farming practices are considered beneficial. For example, the use of shade trees and fallowing all help preserve soil organic matter, which is a critical factor in the maintenance of soil fertility in the deeply weathered and leached soils common in the Amazon.
There is a diversity of forest people in Asia, including the Lumad peoples of the Philippines and the Penan and Dayak people of Borneo. The Dayaks are a particularly interesting group as they are noted for their traditional headhunting culture. Fresh human heads were required to perform certain rituals such as the Iban "kenyalang" and the Kenyah "mamat". Pygmies who live in Southeast Asia are, amongst others, referred to as "Negrito".
Resources
Cultivated foods and spices
Yam, coffee, chocolate, banana, mango, papaya, macadamia, avocado, and sugarcane all originally came from tropical rainforest and are still mostly grown on plantations in regions that were formerly primary forest. In the mid-1980s and 1990s, 40 million tons of bananas were consumed worldwide each year, along with 13 million tons of mango. Central American coffee exports were worth US$3 billion in 1970. Much of the genetic variation used in evading the damage caused by new pests is still derived from resistant wild stock. Tropical forests have supplied 250 cultivated kinds of fruit, compared to only 20 for temperate forests. Forests in New Guinea alone contain 251 tree species with edible fruits, of which only 43 had been established as cultivated crops by 1985.
Ecosystem services
In addition to extractive human uses, rain forests also have non-extractive uses that are frequently summarized as ecosystem services. Rain forests play an important role in maintaining biological diversity, sequestering and storing carbon, global climate regulation, disease control, and pollination. Half of the rainfall in the Amazon area is produced by the forests. The moisture from the forests is important to the rainfall in Brazil, Paraguay, Argentina Deforestation in the Amazon rainforest region was one of the main reason that cause the severe Drought of 2014–2015 in Brazil For the last three decades, the amount of carbon absorbed by the world's intact tropical forests has fallen, according to a study published in 2020 in the journal Nature. In 2019 they took up a third less carbon than they did in the 1990s, due to higher temperatures, droughts and deforestation. The typical tropical forest may become a carbon source by the 2060s.
Tourism
Despite the negative effects of tourism in the tropical rainforests, there are also several important positive effects.
In recent years ecotourism in the tropics has increased. While rainforests are becoming increasingly rare, people are travelling to nations that still have this diverse habitat. Locals are benefiting from the additional income brought in by visitors, as well areas deemed interesting for visitors are often conserved. Ecotourism can be an incentive for conservation, especially when it triggers positive economic change. Ecotourism can include a variety of activities including animal viewing, scenic jungle tours and even viewing cultural sights and native villages. If these practices are performed appropriately this can be beneficial for both locals and the present flora and fauna.
An increase in tourism has increased economic support, allowing more revenue to go into the protection of the habitat. Tourism can contribute directly to the conservation of sensitive areas and habitat. Revenue from park-entrance fees and similar sources can be utilised specifically to pay for the protection and management of environmentally sensitive areas. Revenue from taxation and tourism provides an additional incentive for governments to contribute revenue to the protection of the forest.
Tourism also has the potential to increase public appreciation of the environment and to spread awareness of environmental problems when it brings people into closer contact with the environment. Such increased awareness can induce more environmentally conscious behavior. Tourism has had a positive effect on wildlife preservation and protection efforts, notably in Africa but also in South America, Asia, Australia, and the South Pacific.
Conservation
Threats
Deforestation
Mining and drilling
Deposits of precious metals (gold, silver, coltan) and fossil fuels (oil and natural gas) occur underneath rainforests globally. These resources are important to developing nations and their extraction is often given priority to encourage economic growth. Mining and drilling can require large amounts of land development, directly causing deforestation. In Ghana, a West African nation, deforestation from decades of mining activity left about 12% of the country's original rainforest intact.
Conversion to agricultural land
With the invention of agriculture, humans were able to clear sections of rainforest to produce crops, converting it to open farmland. Such people, however, obtain their food primarily from farm plots cleared from the forest and hunt and forage within the forest to supplement this. The issue arising is between the independent farmer providing for his family and the needs and wants of the globe as a whole. This issue has seen little improvement because no plan has been established for all parties to be aided.
Agriculture on formerly forested land is not without difficulties. Rainforest soils are often thin and leached of many minerals, and the heavy rainfall can quickly leach nutrients from area cleared for cultivation. People such as the Yanomamo of the Amazon, utilize slash-and-burn agriculture to overcome these limitations and enable them to push deep into what were previously rainforest environments. However, these are not rainforest dwellers, rather they are dwellers in cleared farmland that make forays into the rainforest. Up to 90% of the typical Yanamomo diet comes from farmed plants.
Some action has been taken by suggesting fallow periods of the land allowing secondary forest to grow and replenish the soil. Beneficial practices like soil restoration and conservation can benefit the small farmer and allow better production on smaller parcels of land.
Climate change
The tropics take a major role in reducing atmospheric carbon dioxide. The tropics (most notably the Amazon rainforest) are called carbon sinks. As major carbon reducers and carbon and soil methane storages, their destruction contributes to increasing global energy trapping, atmospheric gases. Climate change has been significantly contributed to by the destruction of the rainforests. A simulation was performed in which all rainforest in Africa were removed. The simulation showed an increase in atmospheric temperature by 2.5 to 5 degrees Celsius.
Declining populations
Some species of fauna show a trend towards declining populations in rainforests, for example, reptiles that feed on amphibians and reptiles. This trend requires close monitoring. The seasonality of rainforests affects the reproductive patterns of amphibians, and this in turn can directly affect the species of reptiles that feed on these groups, particularly species with specialized feeding, since these are less likely to use alternative resources.
Protection
Efforts to protect and conserve tropical rainforest habitats are diverse and widespread. Tropical rainforest conservation ranges from strict preservation of habitat to finding sustainable management techniques for people living in tropical rainforests. International policy has also introduced a market incentive program called Reducing Emissions from Deforestation and Forest Degradation (REDD) for companies and governments to outset their carbon emissions through financial investments into rainforest conservation.
| Physical sciences | Forests | null |
931806 | https://en.wikipedia.org/wiki/Greater%20kudu | Greater kudu | The greater kudu (Tragelaphus strepsiceros) is a large woodland antelope, found throughout eastern and southern Africa. Despite occupying such widespread territory, they are sparsely populated in most areas due to declining habitat, deforestation, and poaching. The greater kudu is one of two species commonly known as kudu, the other being the lesser kudu, T. imberbis.
Etymology
Kudu ( ), or koodoo, is the Khoikhoi name for this antelope. Trag- (Greek) denotes a goat and elaphos (Greek) a deer. Strepho (Greek) means 'twist', and strepsis is 'twisting'. Keras (Greek) refers to the horn of the animal.
Physical characteristics
Greater kudus have a narrow body with long legs, and their coats can range from brown/bluish grey to reddish brown. They possess between 4 and 12 vertical white stripes along their torso. The head tends to be darker in colour than the rest of the body, and exhibits a small white chevron which runs between the eyes.
Greater kudu bulls tend to be much larger than the cows, and vocalize much more, utilizing low grunts, clucks, humming, and gasping. The bulls also have beards running along their throats, and large horns with two and a half twists, which, were they to be straightened, would reach an average length of , with the record being . They diverge slightly as they slant back from the head. The horns do not begin to grow until the bull is between the ages of 6–12 months. The horns form the first spiral rotation at around 2 years of age, and not reaching the full two and a half rotations until they are 6 years old; occasionally they may even have 3 full turns.
The greater kudu is one of the largest species of antelope, being slightly smaller than the bongo. Bulls weigh , with a maximum of , and stand up to tall at the shoulder. The ears of the greater kudu are large and round. Cows weigh and stand as little as tall at the shoulder; they are hornless, without a beard or nose markings. The head-and-body length is , to which the tail may add a further .
Taxonomy and subspecies
Formerly four subspecies have been described, but recently only one to three subspecies have been accepted based on colour, number of stripes and horn length:
T. s. strepsiceros – southern parts of the range from southern Kenya to Namibia, Botswana, and South Africa
T. s. chora – northeastern Africa from northern Kenya through Ethiopia to eastern Sudan, Somalia, and Eritrea
T. s. cottoni – Chad and western Sudan
This classification was supported by the genetic difference of one specimen of northern Kenya (T. s. chora) in comparison with several samples from the southern part of the range between Tanzania and Zimbabwe (T. s. strepsiceros). No specimen of the northwestern population, which may represent a third subspecies (T. s. cottoni), was tested within this study.
In Groves and Grubb's book Ungulate Taxonomy, a recent taxonomic revision was made that evaluated all species and subspecies of kudu and other ungulates. This review split the genus Tragelaphus into 4 separate genera, Tragelaphus (bushbuck, sitatunga, bongo, nyala, and gedemsa or mountain nyala), Ammelaphus (lesser kudu), Strepsiceros (greater kudu), and their close relatives Taurotragus (elands). The greater kudu was split into four species based on genetic evidence and morphological features (horn structure and coat color). Each species was based on a different subspecies, Strepsiceros strepsiceros (Cape kudu), Strepsiceros chora (northern kudu), Strepsiceros cottoni (western kudu), and Strepsiceros zambesiensis (Zambezi kudu) which is not commonly accepted even as a subspecies. The Cape kudu is found in south central South Africa, the Zambezi kudu (closely related to the Cape kudu) is found from northern to southern Tanzania and northern South Africa, Namibia, and Angola through Zambia, Mozambique, and eastern DR Congo, the northern kudu is found in eastern Sudan southwards through Ethiopia and Kenya to the Tanzanian border, and the western kudu is found in southeastern Chad, western Sudan, and in northern Central African Republic. Although this alternative taxonomy is not commonly accepted, it was accepted in the Handbook of the Mammals of the World.
Range and ecology
The range of the greater kudu extends from the east in Ethiopia, Tanzania, Eritrea and Kenya into the south where they are found in Zambia, Angola, Namibia, Botswana, Zimbabwe and South Africa. Other regions where greater kudu are located are Central African Republic, Chad, Democratic Republic of the Congo, Djibouti, Eswatini, Malawi, Mozambique, Somalia, and Uganda. They have also been introduced in small numbers into New Mexico, but were never released into the wild. Their habitat includes mixed scrub woodlands (the greater kudu is one of the few largest mammals that prefer living in settled areas – in scrub woodland and bush on abandoned fields and degraded pastures, mopane bush and acacia in lowlands, hills and mountains. They will occasionally venture onto plains only if there is a large abundance of bushes, but normally avoid such open areas to avoid becoming an easy target for their predators. Their diet consists of leaves, grass, shoots and occasionally tubers, roots and fruit (they are especially fond of oranges and tangerines).
During the day, greater kudus normally cease to be active and instead seek cover under woodland, especially during hot days. They feed and drink in the early morning and late afternoon, acquiring water from waterholes or roots and bulbs that have a high water content. Although they tend to stay in one area, the greater kudu may search over a large distance for water in times of drought, in southern Namibia where water is relatively scarce they have been known to cover extensive distances in very short periods of time.
Predation
Predators of the greater kudu generally consist of lions, spotted hyenas, and African wild dogs. Although cheetahs and leopards also prey on greater kudus, they usually target cows and calves rather than fully grown bulls. There are several instances reported where Nile crocodiles have preyed on greater kudus, although based on records, the larger mammalian carnivores statistically are much more dangerous to the kudu and comparable large ungulates, or at least those with a preference for dry, upland habitats over riparian or swamp areas. When a herd is threatened by predators, an adult (usually a female) will issue a bark to alert the rest of the herd. Despite being very nimble over rocky hillsides and mountains, the greater kudu is not fast enough (nor does it have enough endurance) to escape its main predators over open terrain, so it tends to rely on leaping over shrubs and small trees to shake off pursuers. Greater kudus have excellent hearing and acute eyesight, which helps to alert them to approaching predators. Their colouring and markings protect kudus by camouflaging them. If alarmed, they usually stand still, making them very difficult to spot.
Behavior and social organization
Greater kudus have a lifespan of 7 to 8 years in the wild, and up to 23 years in captivity. They may be active throughout the 24-hour day. Herds disperse during the rainy season when food is plentiful. During the dry season, there are only a few concentrated areas of food so the herds will congregate. Greater kudu are not territorial; they have home areas instead. Maternal herds have home ranges of approximately 4 square kilometers and these home ranges can overlap with other maternal herds. Home ranges of adult males are about 11 square kilometers and generally encompass the ranges of two or three female groups. Females usually form small groups of 6–10 with their offspring, but sometimes they can form a herd up to 20 individuals. Male kudus may form small bachelor groups, but they are more commonly found as solitary and widely dispersed individuals. Solitary males will join the group of females and calves (usually 6–10 individuals per group) only during the mating season (April–May in South Africa).
The male kudus are not always physically aggressive with each other, but sparring can sometimes occur between males, especially when both are of similar size and stature. The male kudus exhibit this sparring behavior by interlocking horns and shoving one another. Dominance is established until one male exhibits the lateral display. In rare circumstances, sparring can result in both males being unable to free themselves from the other's horns, which can then result in the death of both animals.
Rarely will a herd reach a size of forty individuals, partly because of the selective nature of their diet which would make foraging for food difficult in large groups. A herd's area can encompass , and spend an average of 54% of the day foraging for food.
Reproduction
Greater kudus reach sexual maturity between ages 1 to 3. The mating season occurs at the end of the rainy season, which can fluctuate slightly according to the region and climate. Before mating, there is a courtship ritual which consists of the male standing in front of the female and often engaging in a neck wrestle. The male then trails the female while issuing a low pitched call until the female allows him to copulate with her. Gestation takes around 240 days (or eight months). Calving generally starts between February and March (late austral summer), when the grass tends to be at its highest.
Greater kudus tend to bear one , although occasionally there may be two. The pregnant female kudu will leave her group to give birth; once she gives birth, the newborn is hidden in vegetation for about 4 to 5 weeks (to avoid predation). After 4 or 5 weeks, the offspring will accompany its mother for short periods of time; then by 3 to 4 months of age, it will accompany her at all times. By the time it is 6 months old, it is quite independent of its mother. The majority of births occur during the wet season (January to March). In terms of maturity, female greater kudus reach sexual maturity at 15–21 months. Males reach maturity at 21–24 months.
Human interaction
Greater kudus have both benefited and suffered from interaction with humans. Humans are turning much of the kudu's natural habitat into farmland, restricting their home ranges. Humans have also destroyed woodland cover, which they use for their habitat. However, wells and irrigation set up by humans has also allowed the greater kudu to occupy territory that would have been too devoid of water for them previously. The greater kudu are also a target for poachers for meat and horns. The horns of greater kudus are commonly used to make Shofars, a Jewish ritual horn blown at Rosh Hashanah. The animal appears on the Eritrean 50-cent coin.
Status
The greater kudu population in the northern part of its range has declined due to excessive hunting and rapid habitat loss. However, they are evaluated as low risk in the IUCN Red List of endangered species. The long-term survival of the greater kudu at large is not in jeopardy as populations located elsewhere remain robust and well-managed.
The greater kudu receives adequate protection from southern Tanzania to South Africa. There are large populations in parks and reserves such as Ruaha-Rungwa-Kisigo and Selous (Tanzania), Luangwa Valley and Kafue (Zambia), Etosha (Namibia), Moremi, Chobe and Central Kalahari (Botswana), Hwange, Chizarira, Mana Pools and Gonarezhou (Zimbabwe) and in Kruger (11,200–17,300) and Hluhluwe–iMfolozi (South Africa). An abundance of greater kudu is also found in private farms and conservancies in southern Africa, in particular in Namibia, Zimbabwe and South Africa, where they are popular amongst trophy hunters.
| Biology and health sciences | Bovidae | Animals |
931973 | https://en.wikipedia.org/wiki/Myiasis | Myiasis | Myiasis ( ), also known as flystrike or fly strike, is the parasitic infestation of the body of a live animal by fly larvae (maggots) that grow inside the host while feeding on its tissue. Although flies are most commonly attracted to open wounds and urine- or feces-soaked fur, some species (including the most common myiatic flies—the botfly, blowfly, and screwfly) can create an infestation even on unbroken skin. Non-myiatic flies (such as the common housefly) can be responsible for accidental myiasis.
Because some animals (particularly non-native domestic animals) cannot react as effectively as humans to the causes and effects of myiasis, such infestations present a severe and continuing problem for livestock industries worldwide, causing severe economic losses where they are not mitigated by human action. Although typically a far greater issue for animals, myiasis is also a relatively frequent disease for humans in rural tropical regions where myiatic flies thrive, and often may require medical attention to surgically remove the parasites.
Myiasis varies widely in the forms it takes and its effects on those affected. Such variations depend largely on the fly species and where the larvae are located. Some flies lay eggs in open wounds, other larvae may invade unbroken skin or enter the body through the nose or ears, and still others may be swallowed if the eggs are deposited on the lips or food. There can also be accidental myiasis that Eristalis tenax can cause in humans via water containing the larvae or in contaminated uncooked food. The name of the condition derives from ancient Greek μυῖα (myia), meaning "fly".
Signs and symptoms
How myiasis affects the human body depends on where the larvae are located. Larvae may infect dead, necrotic (prematurely dying) or living tissue in various sites: the skin, eyes, ears, stomach, and intestinal tract, or in genitourinary sites. They may invade open wounds and lesions or unbroken skin. Some enter the body through the nose or ears. Larvae or eggs can reach the stomach or intestines if they are swallowed with food and cause gastric or intestinal myiasis. In extremely rare cases, maggots may occasionally infest the vulvar area.
Several different presentations of myiasis and their symptoms:
Wound
Wound myiasis occurs when fly larvae infest open wounds. It has been a serious complication of war wounds in tropical areas and is sometimes seen in neglected wounds in most parts of the world. Predisposing factors include poor socioeconomic conditions, extremes of age, neglect, mental disability, psychiatric illness, alcoholism, diabetes, and vascular occlusive disease.
Eye
Myiasis of the human eye or ophthalmomyiasis can be caused by Hypoderma tarandi, a parasitic botfly of caribou. It is known to lead to uveitis, glaucoma, and retinal detachment.
Cause
Life cycle
The life cycle in sheep is typical of the disease. The female flies lay their eggs on the sheep in damp, protected areas of the body that are soaked with urine and feces, mainly the sheep's breech (buttocks). It takes approximately eight hours to a day for the eggs to hatch, depending on the conditions. Once hatched, the larvae then lacerate the skin with their mouthparts, causing open sores. Once the skin has been breached, the larvae then tunnel through the sores into the host's subcutaneous tissue, causing deep and irritating lesions highly subject to infection. After about the second day, bacterial infection is likely and, if left untreated, causes bacterial bloodstream infections or sepsis. This leads to anorexia and weakness and is generally fatal if untreated.
Species affecting humans
There are three main fly families causing economically important myiasis in livestock and also, occasionally, in humans:
Calliphoridae (blowflies)
Some examples include Calliphora vomitoria, Calliphora vicina, and Cordylobia
Oestridae (botflies)
Sarcophagidae (fleshflies) Sarcophaga barbata are usually found in dead and rotting meat and animal excrement, which are prime environments for them. This is because their larvae are facultative parasites, as they feed on organic tissue and use the hosts' oxygen reserve.
Other families occasionally involved are:
Anisopodidae
Piophilidae
Stratiomyidae
Syrphidae
Specific myiasis
Caused by flies that need a host for larval development:
Dermatobia hominis (human botfly)
Cordylobia anthropophaga (tumbu fly)
Cordylobia rodhaini (Lund's fly)
Oestrus ovis (sheep botfly)
Hypoderma spp. (cattle botflies or ox warbles)
Gasterophilus spp. (horse botfly)
Cochliomyia hominivorax (new world screwworm fly)
Chrysomya bezziana (old world screwworm fly)
Auchmeromyia senegalensis (Congo floor maggot)
Cuterebra spp. (rodent and rabbit botfly)
Semispecific myiasis
Caused by flies that usually lay their eggs in decaying animal or vegetable matter, but that can develop in a host if open wounds or sores are present:
Lucilia spp. (green-bottle fly)
Cochliomyia spp. (screw-worm fly)
Phormia spp. (black-bottle fly)
Calliphora spp. (blue-bottle fly)
Sarcophaga spp. (flesh fly or sarcophagids)
Flesh flies, or sarcophagids, members of the family Sarcophagidae, can cause intestinal myiasis in humans if the females lay their eggs on meat or fruit.
Accidental myiasis
Accidental myiasis is also called pseudomyiasis. It is caused by flies that have no preference or need to develop in a host but may do so on rare occasions. Transmission occurs through accidental deposit of eggs on oral or genitourinary openings, or by swallowing eggs or larvae that are on food. The cheese fly (Piophila casei) sometimes causes myiasis through intentional consumption of its maggots (which are contained in the traditional Sardinian delicacy casu marzu). Other flies that can accidentally cause myiasis are:
Musca domestica (housefly)
Fannia spp. (latrine flies)
Eristalis tenax (rat-tailed maggots)
Muscina spp.
The adult flies are not parasitic, but when they lay their eggs in open wounds and these hatch into their larval stage (also known as maggots or grubs), the larvae feed on live or necrotic tissue, causing myiasis to develop. They may also be ingested or enter through other body apertures.
Diagnosis
Myiasis is often misdiagnosed in the United States because it is rare and its symptoms are not specific. Intestinal myiasis and urinary myiasis are especially difficult to diagnose.
Clues that myiasis may be present include recent travel to an endemic area, one or more non-healing lesions on the skin, itchiness, movement under the skin or pain, discharge from a central punctum (tiny hole), or a small, white structure protruding from the lesion. Serologic testing has also been used to diagnose the presence of botfly larvae in human ophthalmomyiasis.
Classifications
German entomologist Fritz Zumpt describes myiasis as "the infestation of live human and vertebrate animals with dipterous larvae, which at least for a period, feed on the host's dead or living tissue, liquid body substances, or ingested food". For modern purposes, however, this is too vague. For example, feeding on dead or necrotic tissue is not generally a problem except when larvae such as those of flies in the family Piophilidae attack stored food such as cheese or preserved meats; such activity suggests saprophagy rather than parasitism; it even may be medically beneficial in maggot debridement therapy (MDT).
Currently, myiasis commonly is classified according to aspects relevant to the case in question:
The classical description of myiasis is according to the part of the host that is infected. This is the classification used by ICD-10. For example:
dermal
sub-dermal
cutaneous (B87.0)
creeping, where larvae burrow through or under the skin
furuncular, where a larva remains in one spot, causing a boil-like lesion
nasopharyngeal, in the nose, sinuses or pharynx (B87.3)
ophthalmic or ocular, in or about the eye (B87.2)
auricular, in or about the ear
gastric, rectal, or intestinal/enteric for the appropriate part of the digestive system (B87.8)
urogenital (B87.8)
Another aspect is the relationship between the host and the parasite which provides insight into the biology of the fly species causing the myiasis and its likely effect. Thus the myiasis is described as either:
obligatory, where the parasite cannot complete its life cycle without its parasitic phase, which may be specific, semispecific, or opportunistic
facultative, incidental, or accidental, where it is not essential to the life cycle of the parasite; perhaps a normally free-living larva accidentally gained entrance to the host
Accidental myiasis commonly is enteric, resulting from swallowing eggs or larvae with one's food. The effect is called pseudomyiasis. One traditional cause of pseudomyiasis was the eating of maggots of cheese flies in cheeses such as Stilton. Depending on the species present in the gut, pseudomyiasis may cause significant medical symptoms, but it is likely that most cases pass unnoticed.
Prevention
The first control method is preventive and aims to eradicate the adult flies before they can cause any damage and is called. The second control method is treatment once the infestation is present, and concerns the infected animals (including humans).
The principal control method of adult populations of myiasis-inducing flies involves insecticide applications in the environment where the target livestock is kept. Organophosphorus or organochlorine compounds may be used, usually in a spraying formulation. One alternative prevention method is the sterile insect technique (SIT) where a significant number of artificially reared sterilized (usually through irradiation) male flies are introduced. The male flies compete with wild breed males for females to copulate and thus cause females to lay batches of unfertilized eggs that cannot develop into the larval stage.
One prevention method involves removing the environment most favourable to the flies, such as by removal of the tail. Another example is the crutching of sheep, which involves the removal of wool from around the tail and between the rear legs, which is a favourable environment for the larvae. Another, more permanent, practice that is used in some countries is mulesing, where the skin is removed from young animals to tighten remaining skinleaving it less prone to fly attack.
To prevent myiasis in humans, there is a need for general improvement of sanitation, personal hygiene, and extermination of the flies by insecticides. Clothes should be washed thoroughly, preferably in hot water, dried away from flies, and ironed thoroughly. The heat of the iron kills the eggs of myiasis-causing flies.
Treatment
This applies once an infestation is established. In many circles the first response to cutaneous myiasis once the breathing hole has formed, is to cover the air hole thickly with petroleum jelly. Lack of oxygen then forces the larva to the surface, where it can more easily be dealt with. In a clinical or veterinary setting there may not be time for such tentative approaches, and the treatment of choice might be more direct, with or without an incision. First, the larva must be eliminated through pressure around the lesion and the use of forceps. Secondly, the wound must be cleaned and disinfected. Further control is necessary to avoid further reinfestation.
Livestock may be treated prophylactically with slow-release boluses containing ivermectin, which can provide long-term protection against the development of the larvae. Sheep also may be dipped, a process that involves drenching the animals in persistent insecticide to poison the larvae before they develop into a problem.
Epidemiology
Myiasis is prevalent in livestock, and especially in domestic sheep. Myiasis in sheep is often caused by blowflies (Lucilia sericata and L. cuprina in particular) and is commonly referred to as blowfly strike. Blowfly strike, and other flystrike, occurs worldwide but is most common in regions where hot and wet conditions are sustained, such as Sub-Saharan Africa, Southeast Asia, Latin America, Australia, and New Zealand. As of 2021, blowfly strike accounts for over A$280 million a year in losses for the Australian sheep industry. As mitigation, Australian sheep farmers may engage in mulesing, a procedure designed to remove strips of wool-producing skin that are the most common targets for flies. Farmers may also dock lambs' tails to reduce the likelihood of infestation. However, both mulesing and tail-docking have received criticism from animal welfare groups, who say the mitigative procedures are excessive and can have other negative effects.
In addition to blowfly strike in sheep, myiasis from screwworm flies (Cochliomyia hominivorax in particular) regularly cause upwards of US$100 million in annual damages to domestic cows and goats. Screwworm-related myiasis is primarily mitigated through the sterile insect technique.
History
Frederick William Hope coined the term myiasis in 1840 to refer to diseases resulting from dipterous larvae as opposed to those caused by other insect larvae (the term for this was scholechiasis). Hope described several cases of myiasis from Jamaica caused by unknown larvae, one of which resulted in death.
Even though the term myiasis was first used in 1840, such conditions have been known since ancient times. Ambroise Paré, the chief surgeon to King Charles IX and King Henry III, observed that maggots often infested open wounds.
Maggot therapy
Throughout recorded history, maggots have been used therapeutically to clean out necrotic wounds, an application known as maggot therapy''.
Fly larvae that feed on dead tissue can clean wounds and may reduce bacterial activity and the chance of a secondary infection. They dissolve dead tissue by secreting digestive enzymes onto the wound as well as actively eating the dead tissue with mouth hooks, two hard, probing appendages protruding on either side of the "mouth". Maggot therapyalso known as maggot debridement therapy (MDT), larval therapy, larva therapy, or larvae therapyis the intentional introduction by a health care practitioner of live, disinfected green bottle fly maggots into the non-healing skin and soft tissue wounds of a human or other animal for the purpose of selectively cleaning out only the necrotic tissue within a wound to promote healing.
Although maggot therapy has been used in the US for the past 80 years, it was approved by the FDA as a medical device only in 2004 (along with leeches). Maggots were the first live organism to be marketed in the US according to FDA regulations, and are approved for treating neuropathic (diabetic) foot ulcers, pressure ulcers, venous stasis ulcers, and traumatic and post-surgical wounds that are unresponsive to conventional therapies. Maggots were used in medicine before this time, but were not federally regulated. In 1990, California internist Ronald Sherman began treating patients with maggots produced at his lab at the UC Irvine School of Medicine. Sherman went on to co-found Monarch Labs in 2005, which UC Irvine contracted to produce maggots for Sherman's own continuing clinical research on myiasis at the university. Monarch Labs also sells maggots to hospitals and other medical practices, the first US commercial supplier to do so since the last one closed in 1935.
In the US, demand for these fly larvae doubled after the FDA ruling. Maggot therapy is now used in more than 300 sites across the country. The American Medical Association and Centers for Medicare and Medicaid Services recently clarified the reimbursement guidelines to the wound care community for medicinal maggots, and this therapy may soon be covered by insurance. The larvae of the green bottle fly (Lucilia fly) are now used exclusively for this purpose, since they preferentially devour only necrotic tissue, leaving healthy tissue intact. This is an important distinction, as most other major varieties of myiasitic fly larvae attack both live and dead wound tissue indiscriminately, effectively negating their benefit in non-harmful wound debridement. Medicinal maggots are placed on the wound and covered with a sterile dressing of gauze and nylon mesh. However, too many larvae placed on the wound could result in healthy tissue being eaten, efficiently creating a new wound, and rendering it a type of myiasis.
History
Maggot therapy has a long history and prehistory. The indigenous people of Australia used maggot therapy, and so do the Hill Peoples of Northern Burma, and possibly the Mayans of Central America. Surgeons in Napoleon's armies recognized that wounded soldiers with myiasis were more likely to survive than those without the infestation. In the American Civil War, army surgeons treated wounds by allowing blowfly maggots to clean away the decayed tissue.
William Baer, an orthopedic surgeon at Johns Hopkins during the late 1920s, used maggot therapy to treat a series of patients with osteomyelitis, an infection of bone or bone marrow. The idea was based on an experience in World War I in which two soldiers presented to him with broken femurs after having lain on the ground for seven days without food. Baer could not figure out why neither man had a fever or signs of sepsis. He observed: "On removing the clothing from the wounded part, much was my surprise to see the wound filled with thousands and thousands of maggots, apparently those of the blow fly. The sight was very disgusting and measures were taken hurriedly to wash out these abominable looking creatures." However, he then saw that the wounds were filled with "beautiful pink granulation tissue" and were healing well.
Maggot therapy was common in the United States during the 1930s. However, during the second half of the twentieth century, after the introduction of antibiotics, maggot therapy was used only as a last resort for very serious wounds. Lately maggots have been making a comeback due to the increased resistance of bacteria to antibiotics.
| Biology and health sciences | Helminthic diseases and infestations | Health |
932217 | https://en.wikipedia.org/wiki/Classical%20unified%20field%20theories | Classical unified field theories | Since the 19th century, some physicists, notably Albert Einstein, have attempted to develop a single theoretical framework that can account for all the fundamental forces of nature – a unified field theory. Classical unified field theories are attempts to create a unified field theory based on classical physics. In particular, unification of gravitation and electromagnetism was actively pursued by several physicists and mathematicians in the years between the two World Wars. This work spurred the purely mathematical development of differential geometry.
This article describes various attempts at formulating a classical (non-quantum), relativistic unified field theory. For a survey of classical relativistic field theories of gravitation that have been motivated by theoretical concerns other than unification, see Classical theories of gravitation. For a survey of current work toward creating a quantum theory of gravitation, see quantum gravity.
Overview
The early attempts at creating a unified field theory began with the Riemannian geometry of general relativity, and attempted to incorporate electromagnetic fields into a more general geometry, since ordinary Riemannian geometry seemed incapable of expressing the properties of the electromagnetic field. Einstein was not alone in his attempts to unify electromagnetism and gravity; a large number of mathematicians and physicists, including Hermann Weyl, Arthur Eddington, and Theodor Kaluza also attempted to develop approaches that could unify these interactions. These scientists pursued several avenues of generalization, including extending the foundations of geometry and adding an extra spatial dimension.
Early work
The first attempts to provide a unified theory were by G. Mie in 1912 and Ernst Reichenbacher in 1916. However, these theories were unsatisfactory, as they did not incorporate general relativity because general relativity had yet to be formulated. These efforts, along with those of Rudolf Förster, involved making the metric tensor (which had previously been assumed to be symmetric and real-valued) into an asymmetric and/or complex-valued tensor, and they also attempted to create a field theory for matter as well.
Differential geometry and field theory
From 1918 until 1923, there were three distinct approaches to field theory: the gauge theory of Weyl, Kaluza's five-dimensional theory, and Eddington's development of affine geometry. Einstein corresponded with these researchers, and collaborated with Kaluza, but was not yet fully involved in the unification effort.
Weyl's infinitesimal geometry
In order to include electromagnetism into the geometry of general relativity, Hermann Weyl worked to generalize the Riemannian geometry upon which general relativity is based. His idea was to create a more general infinitesimal geometry. He noted that in addition to a metric field there could be additional degrees of freedom along a path between two points in a manifold, and he tried to exploit this by introducing a basic method for comparison of local size measures along such a path, in terms of a gauge field. This geometry generalized Riemannian geometry in that there was a vector field Q, in addition to the metric g, which together gave rise to both the electromagnetic and gravitational fields. This theory was mathematically sound, albeit complicated, resulting in difficult and high-order field equations. The critical mathematical ingredients in this theory, the Lagrangians and curvature tensor, were worked out by Weyl and colleagues. Then Weyl carried out an extensive correspondence with Einstein and others as to its physical validity, and the theory was ultimately found to be physically unreasonable. However, Weyl's principle of gauge invariance was later applied in a modified form to quantum field theory.
Kaluza's fifth dimension
Kaluza's approach to unification was to embed space-time into a five-dimensional cylindrical world, consisting of four space dimensions and one time dimension. Unlike Weyl's approach, Riemannian geometry was maintained, and the extra dimension allowed for the incorporation of the electromagnetic field vector into the geometry. Despite the relative mathematical elegance of this approach, in collaboration with Einstein and Einstein's aide Grommer it was determined that this theory did not admit a non-singular, static, spherically symmetric solution. This theory did have some influence on Einstein's later work and was further developed later by Klein in an attempt to incorporate relativity into quantum theory, in what is now known as Kaluza–Klein theory.
Eddington's affine geometry
Sir Arthur Stanley Eddington was a noted astronomer who became an enthusiastic and influential promoter of Einstein's general theory of relativity. He was among the first to propose an extension of the gravitational theory based on the affine connection as the fundamental structure field rather than the metric tensor which was the original focus of general relativity. Affine connection is the basis for parallel transport of vectors from one space-time point to another; Eddington assumed the affine connection to be symmetric in its covariant indices, because it seemed plausible that the result of parallel-transporting one infinitesimal vector along another should produce the same result as transporting the second along the first. (Later workers revisited this assumption.)
Eddington emphasized what he considered to be epistemological considerations; for example, he thought that the cosmological constant version of the general-relativistic field equation expressed the property that the universe was "self-gauging". Since the simplest cosmological model (the De Sitter universe) that solves that equation is a spherically symmetric, stationary, closed universe (exhibiting a cosmological red shift, which is more conventionally interpreted as due to expansion), it seemed to explain the overall form of the universe.
Like many other classical unified field theorists, Eddington considered that in the Einstein field equations for general relativity the stress–energy tensor , which represents matter/energy, was merely provisional, and that in a truly unified theory the source term would automatically arise as some aspect of the free-space field equations. He also shared the hope that an improved fundamental theory would explain why the two elementary particles then known (proton and electron) have quite different masses.
The Dirac equation for the relativistic quantum electron caused Eddington to rethink his previous conviction that fundamental physical theory had to be based on tensors. He subsequently devoted his efforts into development of a "Fundamental Theory" based largely on algebraic notions (which he called "E-frames"). Unfortunately his descriptions of this theory were sketchy and difficult to understand, so very few physicists followed up on his work.
Einstein's geometric approaches
When the equivalent of Maxwell's equations for electromagnetism is formulated within the framework of Einstein's theory of general relativity, the electromagnetic field energy (being equivalent to mass as defined by Einstein's equation E=mc2) contributes to the stress tensor and thus to the curvature of space-time, which is the general-relativistic representation of the gravitational field; or putting it another way, certain configurations of curved space-time incorporate effects of an electromagnetic field. This suggests that a purely geometric theory ought to treat these two fields as different aspects of the same basic phenomenon. However, ordinary Riemannian geometry is unable to describe the properties of the electromagnetic field as a purely geometric phenomenon.
Einstein tried to form a generalized theory of gravitation that would unify the gravitational and electromagnetic forces (and perhaps others), guided by a belief in a single origin for the entire set of physical laws. These attempts initially concentrated on additional geometric notions such as vierbeins and "distant parallelism", but eventually centered around treating both the metric tensor and the affine connection as fundamental fields. (Because they are not independent, the metric-affine theory was somewhat complicated.) In general relativity, these fields are symmetric (in the matrix sense), but since antisymmetry seemed essential for electromagnetism, the symmetry requirement was relaxed for one or both fields. Einstein's proposed unified-field equations (fundamental laws of physics) were generally derived from a variational principle expressed in terms of the Riemann curvature tensor for the presumed space-time manifold.
In field theories of this kind, particles appear as limited regions in space-time in which the field strength or the energy density is particularly high. Einstein and coworker Leopold Infeld managed to demonstrate that, in Einstein's final theory of the unified field, true singularities of the field did have trajectories resembling point particles. However, singularities are places where the equations break down, and Einstein believed that in an ultimate theory the laws should apply everywhere, with particles being soliton-like solutions to the (highly nonlinear) field equations. Further, the large-scale topology of the universe should impose restrictions on the solutions, such as quantization or discrete symmetries.
The degree of abstraction, combined with a relative lack of good mathematical tools for analyzing nonlinear equation systems, make it hard to connect such theories with the physical phenomena that they might describe. For example, it has been suggested that the torsion (antisymmetric part of the affine connection) might be related to isospin rather than electromagnetism; this is related to a discrete (or "internal") symmetry known to Einstein as "displacement field duality".
Einstein became increasingly isolated in his research on a generalized theory of gravitation, and most physicists consider his attempts ultimately unsuccessful. In particular, his pursuit of a unification of the fundamental forces ignored developments in quantum physics (and vice versa), most notably the discovery of the strong nuclear force and weak nuclear force.
Schrödinger's pure-affine theory
Inspired by Einstein's approach to a unified field theory and Eddington's idea of the affine connection as the sole basis for differential geometric structure for space-time, Erwin Schrödinger from 1940 to 1951 thoroughly investigated pure-affine formulations of generalized gravitational theory. Although he initially assumed a symmetric affine connection, like Einstein he later considered the nonsymmetric field.
Schrödinger's most striking discovery during this work was that the metric tensor was induced upon the manifold via a simple construction from the Riemann curvature tensor, which was in turn formed entirely from the affine connection. Further, taking this approach with the simplest feasible basis for the variational principle resulted in a field equation having the form of Einstein's general-relativistic field equation with a cosmological term arising automatically.
Skepticism from Einstein and published criticisms from other physicists discouraged Schrödinger, and his work in this area has been largely ignored.
Later work
After the 1930s, progressively fewer scientists worked on classical unification, due to the continued development of quantum-theoretical descriptions of the non-gravitational fundamental forces of nature and the difficulties encountered in developing a quantum theory of gravity. Einstein pressed on with his attempts to theoretically unify gravity and electromagnetism, but he became increasingly isolated in this research, which he pursued until his death. Einstein's celebrity status brought much attention to his final quest, which ultimately saw limited success.
Most physicists, on the other hand, eventually abandoned classical unified theories. Current mainstream research on unified field theories focuses on the problem of creating a quantum theory of gravity and unifying with the other fundamental theories in physics, all of which are quantum field theories. (Some programs, such as string theory, attempt to solve both of these problems at once.) Of the four known fundamental forces, gravity remains the one force for which unification with the others proves problematic.
Although new "classical" unified field theories continue to be proposed from time to time, often involving non-traditional elements such as spinors or relating gravitation to an electromagnetic force, none have been generally accepted by physicists yet.
| Physical sciences | Particle physics: General | Physics |
933173 | https://en.wikipedia.org/wiki/Archelon | Archelon | Archelon is an extinct marine turtle from the Late Cretaceous, and is the largest turtle ever to have been documented, with the biggest specimen measuring from head to tail and in body mass. It is known only from the Pierre Shale and has one species, A. ischyros. In the past, the genus also contained A. marshii and A. copei, though these have been reassigned to Protostega and Kansastega, respectively. The genus was named in 1895 by American paleontologist George Reber Wieland based on a skeleton from South Dakota, who placed it into the extinct family Protostegidae. The leatherback sea turtle (Dermochelys coriacea) was once thought to be its closest living relative, but now, Protostegidae is thought to be a completely separate lineage from any living sea turtle.
Archelon had a leathery carapace instead of the hard shell seen in most sea turtles. The carapace may have featured a row of small ridges, each peaking at in height. It had an especially hooked beak and its jaws were adept at crushing, so it probably ate hard-shelled crustaceans, mollusks, and possibly even sponges, while slowly moving over the seafloor. It also potentially consumed other animals, whilst swimming closer to the surface, like jellyfish, squid, or nautiloids. However, its beak may have been better-adapted for shearing flesh, with fish being another possible prey choice. With its large and strong foreflippers, Archelon was likely able to produce powerful strokes necessary for open-ocean travel and, if need be, escape from fellow marine predators. It inhabited the northern Western Interior Seaway, a mild to cool temperate area, dominated by plesiosaurs, hesperornithiform seabirds, and mosasaurs. It may have gone extinct due to the shrinking of the seaway, increased infant mortality rates (in the sea), higher instances of egg and hatchling predation (on land), and a rapidly cooling climate.
Research history
The holotype specimen, YPM 3000, was collected from the Late Campanian-age Pierre Shale of South Dakota along the Cheyenne River in Custer County by American paleontologist George Reber Wieland in 1895, and described by him the following year based on a mostly complete skeleton excluding the skull. He named it Archelon ischyros, genus name from the Ancient Greek - (-) 'first/early', () 'turtle', and species name from () 'mighty' or 'powerful'. A second specimen, a skull, was discovered in 1897 in the same region.
In 1900, Wieland described a second species, A. marshii, from remains collected in 1898 by American paleontologist Othniel Charles Marsh, to whom the species name refers, on the basis that the shell underside (plastron) was thicker and the humeri were straighter. However, in 1909, Wieland reclassified it as Protostega marshii. In 1902, a third, mostly complete specimen was collected also along the Cheyenne River. In the same study, Protostega copei from Kansas, which was first described by Wieland in 1909 and named in honor of Edward Drinker Cope who first erected the family Protostegidae, was moved to the genus Archelon as A. copei. In 1998, A. copei was moved to the new genus (originally named Microstega, but subsequently renamed Kansastega) as K. copei. In 1992, a fourth and the largest specimen to date, nicknamed "Brigitta", was discovered in Oglala Lakota County, South Dakota and resides in the Natural History Museum Vienna. In 2002, a fifth specimen, a partial skeleton, was discovered from the Pierre Shale of North Dakota along the Sheyenne River near Cooperstown.
Description
The holotype measures from head to tail, with the head measuring , the neck , the thoracic vertebrae , the sacrum , and the tail . The largest specimen, Brigitta, measures around from head to tail and from flipper to flipper, and, in life, weighed around . Skulls of Archelon measured at up to .
Archelon had a distinctly elongated and narrow head. It had a defined hooked beak which was probably covered in a sheath in life, reminiscent of the beaks of birds of prey. However, in the back, the cutting edge of the beak is dull compared to such animals. Much of the length of the head derives from the elongated premaxillae–which is the front part of the beak in this animal–and maxillae. The jugal bones, the cheek bones, due to the elongate head, do not project as far as they do in other turtles. The nostrils are elongated and rest on the top of the skull, slightly posited forward, and are unusually horizontal compared to sea turtles. The jugal bones (cheekbones) are rounded as opposed to triangular in sea turtles. The articular bone, which formed the jaw joint, was probably heavily encased in cartilage. The jaw probably moved in a hammering motion.
Five neck vertebrae were recovered from the holotype, and it probably had eight in total in life; they are X-shaped, procoelous–concave on the side towards the head and convex on the other–and their thick frame indicates strong neck muscles. Ten thoracic vertebrae were found, increasing in size until the sixth then rapidly decreasing, and they have little connection with the carapace. The three vertebrae of the sacrum are short and flat. It probably had eighteen tail vertebrae; the first eight to ten (probably in the same area as the carapace) had neural arches, whereas the remaining did not. Its tail likely had a wide range of mobility, and the tail is thought to have been able to bend at nearly a 90° angle horizontally.
The humeri in the upper arms are proportionally massive, and the radii and ulnae of the forearms are short and compact, indicating the animal had strong flippers in life. The flippers would have had a spread of between , though most likely the more conservative estimate. Stretch marks on the limb bones indicate fast growth, with similarities to the leatherback sea turtle, the fastest growing turtle known, whose juveniles have an average growth rate of per year.
Carapace
The carapace comprises on either side eight neuralia–the plates closest to the midline–and nine pleuralia–the plates that connect the midline to the ribs. The plates of the carapace are mostly uniform in dimensions, with the exception of the two pairs of plates corresponding to the eighth thoracic vertebra which are smaller than the others, and the pygal plate closest to the tail which is larger. Archelon has ten pairs of ribs, and, like the leatherback sea turtle but unlike other sea turtles, the first rib does not meet the first pleural. As in sea turtles, the first rib is noticeably shorter than the second, in this case, three quarters of the length. The second to fifth ribs project at a right angle from the midline, and, in the holotype, each measure in length. A rib increases in thickness in the vertical direction distally, as it gets farther from the midline, and the ribs are relatively larger and more well-developed than those of sea turtles. The second to fifth ribs, in the holotype, originate with a thickness of and terminate with around in thickness.
The neuralia and pleuralia form highly irregular and finger-like sutures where they meet, and one plate may have lain over the other plate while the bone was still developing and malleable. The neuralia and pleuralia–the bony portions of the carapace–are particularly thin, and the ribs, especially the first rib, and the shoulder girdle are unusually heavy and may have had to carry extra stress to compensate, a condition seen in ancient ancestral turtles. Archelon has osteosclerotic structures, where the bone is dense and heavy, which probably served as ballasts in life similar to the limb bones of whales and other open-ocean animals.
The carapace, in life, probably featured a row of ridges along the midline over the chest region, perhaps totaling in seven ridges, with each ridge peaking at either . In the absence of firmly jointed neck and pleural plates, the skin over the carapace was probably thick, strong, and leathery in order to compensate and properly support the shoulder girdle. This leathery carapace is also seen in the leatherback sea turtle. The spongy makeup is similar to the bones seen in open-ocean going vertebrates such as dolphins or ichthyosaurs, and was probably also an adaptation to reduce overall weight.
Plastron
A turtle plastron, its underside, comprises (from head-most to tail-most) the epiplastron, the entoplastron, which is small and wedged in between the former and the hyoplastron; then is the hypoplastron and finally, the xiphiplastron. The plastron, as a whole, is thick, and measures (in a specimen described in 1898) in length. Unlike the carapace, it features striations throughout.
In protostegids, the epiplastron and entoplastron are fused together, forming a single unit called an "entepiplastron" or a "paraplastron." This entepiplastron is T-shaped, as opposed to the Y-shaped entoplastrons in other turtles. The top edge of the T rounds off, except at the center which features a small projection. The outward side is slightly convex and bends somewhat, away from the body. The two ends of the T flatten out, getting broader and thinner, as they get farther from the center.
A thick, continuous ridge connects the hyoplastron, hypoplastron, and xiphiplastron. The hyoplastron features a large number of spines projecting around the circumference. The hyoplastron is slightly elliptical, and grows thinner as it gets farther from the center, before the spines erupt. The spines grow thick and narrow towards their middle portion. The seven to nine spines projecting towards the head are short and triangular. The six middle spines are long and thin. The last 19 spines are flat. There are no marks indicating contact with the entepiplastron. The hypoplastron is similar to the hyoplastron, except it has more spines, a total of 54. The xiphiplastron is boomerang-shaped, a primitive characteristic in contrast to the straight ones seen in more modern turtles.
Classification and evolution
In its original 1896 description, Wieland placed Archelon into the family Protostegidae, which included at the time the smaller Protostega and Protosphargis, the latter being nowadays classified in a position closer to the family Cheloniidae. In 1953, Swiss paleontologist Rainer Zangerl split Protostegidae into two subfamilies: Chelospharginae and Protosteginae; to the former was assigned Chelosphargis and Calcarichelys, and the latter Archelon and Protostega. The sister group of Protostegidae had, in the past, been considered to be Dermochelyidae, and thus their closest living relative would have been the dermochelyid leatherback sea turtle (Dermochelys coriacea). However, phylogenetic studies conclude that protostegids represent a completely separate, basal (ancient) lineage that originated in the Late Jurassic, removing the family from the superfamily Chelonioidea (which includes all sea turtles). In this model, Archelon does not share a marine ancestor with any sea turtle.
Paleobiology
Archelon was an obligate carnivore. The thick plastron indicates the animal probably spent a lot of time on the soft, muddy seafloor, likely a slow-moving bottom feeder. According to American paleontologist Samuel Wendell Williston, the jaws were adapted for crushing, implying the turtle ate large mollusks and crustaceans. In 1914, he suggested that the abundant, thin-shelled, bottom-dwelling Cretaceous bivalves–some exceeding in diameter–would have easily been able to sustain Archelon. However, these were probably absent in the central Western Interior Seaway by the Early Campanian. Conversely, the beak may have been adapted for shearing flesh. It might have been able to target larger fish and reptiles, as well as, similar to the leatherback sea turtle, soft-bodied creatures such as squid and jellyfish. However, it is possible the sharp beak was used only in combat against other Archelon. The nautilus Eutrephoceras dekayi was found in great number near an Archelon specimen, and may have been a potential food source. Archelon may have also occasionally scavenged off the surface water.
Archelon probably had weaker arms, and thus less swimming power, than the leatherback sea turtle, and so did not frequent the open ocean as much, preferring shallower, calmer waters. This is indicated by the similarity of the humerus/arm and hand/arm ratios of it and cheloniids, which are known to have poor development of the limbs into flippers and a preference for shallow water. Conversely, the large flipper-to-carapace ratio of protostegids and the similarly large flipper spread, like that of the predatory cheloniid loggerhead sea turtle (Caretta caretta), combined with a broad body, indicate they could have pursued active prey, though they probably could not have sustained high speeds. Overall, it may have been a moderately-good swimmer, capable of open-ocean travel.
Archelon, like other marine turtles, probably would have had to have come onshore to nest; like other turtles, it probably dug out a hollow in the sand, laid several dozens of eggs, and took no part in child rearing. The right lower flipper of the holotype is missing, and stunted growth of remaining flipper bone indicates this occurred early in life. It may have been the result of attempted predation by a bird while a hatchling and trying to escape to the sea, bitten off by some large predator such as a mosasaur or a Xiphactinus, or was crushed off by larger adults while herding on the shore. However, the latter is unlikely as juveniles probably did not frequent the coasts even during breeding season. Brigitta is estimated to have lived to 100 years, and may have died while partially covered in mud brumating–a state of dormancy–on the ocean floor. However, the longstanding belief that marine turtles brumate underwater like freshwater turtles may be incorrect given the high surfacing-frequency needed to prevent drowning.
Paleoecology
Archelon inhabited the shallow Western Interior Seaway; the muddy, oxygen-depleted seafloor was probably, on average, no more than below the surface, and average water temperature may have been in the Campanian. The Late Cretaceous Dakotas were submerged in the Northern Inland Subprovince, an area characterized by moderate to cool temperatures, with an abundance of plesiosaurs, hesperornithiform seabirds, and mosasaurs, particularly Platecarpus. There is no fossil evidence for vertebrate migration between the northern and southern provinces. Though sharks were generally more common in the southern province, several sharks are known from the Pierre Shale, including Squalus, Squalicorax, Pseudocorax, and Cretolamna. Other large predatory fish include ichthyodectids such as Xiphactinus. The Pierre Shale's invertebrate assemblage includes a variety of mollusks, namely ammonites–from the Pierre Shale Placenticeras placenta, Scaphites nodosus, Didymoceras, and Baculites ovatus–bivalves–such as the giant Inoceramus– the squid-like belemnites, and nautilus.
As the seaway progressively migrated southward, it is possible Archelon was unable to migrate with it. The increasing threat of egg or hatchling predation by new marine or mammalian species may have led to the extinction of Archelon, and the disappearance of gigantic protostegids seems to have coincided with the increasing size of dermochelyids. Protostegidae is more-or-less absent in Maastrichtian deposits, the latest Cretaceous, and probably died off due to a cooling trend which other turtles were able to survive due to some thermoregulatory capabilities. Average water temperature may have decreased to depending on estimated CO2 levels. However, some Maastrichtian-age Kansas Pierre Shale fossils may have been eroded millions of years ago, and it is possible Archelon survived well into the Maastrichtian.
| Biology and health sciences | Prehistoric turtles | Animals |
934317 | https://en.wikipedia.org/wiki/Steppe%20bison | Steppe bison | The steppe bison or steppe wisent (Bison priscus) is an extinct species of bison. It was widely distributed across the mammoth steppe, ranging from Western Europe to eastern Beringia in North America during the Late Pleistocene. It is ancestral to all North American bison, including ultimately modern American bison. Three chronological subspecies, Bison priscus priscus, Bison priscus mediator, and Bison priscus gigas, have been suggested.
Evolution
The steppe bison first appeared during the mid Middle Pleistocene in eastern Eurasia, subsequently dispersing westwards as far as Western Europe. During the late Middle Pleistocene, around 195,000-135,000 years ago, the steppe bison migrated across the Bering land bridge into North America, becoming ancestral to endemic North American bison species, including the largest known bison, the long-horned Bison latifrons, and the smaller Bison antiquus, the latter of which is thought to be ancestral to modern American bison.
Description
Resembling the modern bison species, especially the American wood bison (Bison bison athabascae), the steppe bison was over tall at the withers, reaching in weight. The tips of the horns were a meter apart, the horns themselves being over half a meter long.
Bison priscus gigas is the largest known bison of Eurasia. This subspecies was possibly analogous to Bison latifrons, attaining similar body sizes and horns which were up to apart, and presumably favored similar habitat conditions.
The steppe bison was also anatomically similar to the European bison (Bison bonasus), to the point of difficulty distinguishing between the two when complete skeletons are unavailable. The two species were close enough to interbreed; however they were also genetically distinct, indicating that interbreeding was in fact rare, possibly as a result of niche partitioning between the species.
Palaeoecology
Dental microwear analysis suggests the steppe bison was a mixed feeder, rather than a strict grazer. Likely predators of steppe bison include cave hyenas (with steppe bison remains having been found in their dens), cave lions (claw marks were found on the “Blue Babe” mummy and were found to be made by one), as well as scimitar toothed cats and possibly wolves.
Extinction
The steppe bison distribution contracted to the north after the end of the last glacial period, surviving into the mid Holocene before becoming extinct as part of the Quaternary extinction event. A steppe bison skeleton was radiocarbon dated to 5,400 years Before Present (c. 3450 BCE) in Alaska. B. priscus remains in the northern Angara River in Asia were dated to 2550-2450 BCE, and in the Oyat River in Leningrad Oblast, Russia to 1130-1060 BCE. The causes for the extinction of the steppe bison and many other primarily megafaunal species remain hotly debated, but the selectivity for large animals suggests that the spread of modern humans played a substantial role.
Discoveries
Steppe bison appear in cave art, notably in the Cave of Altamira and Lascaux, and the carving Bison Licking Insect Bite, and have been found in naturally ice-preserved form.
Blue Babe is the 36,000-year-old mummy of a male steppe bison which was discovered north of Fairbanks, Alaska, in July 1979. The mummy was noticed by a gold miner who named the mummy Blue Babe – "Babe" for Paul Bunyan's mythical giant ox, permanently turned blue when he was buried to the horns in a blizzard (Blue Babe's own bluish cast was caused by a coating of vivianite, a blue iron phosphate covering much of the specimen). Claw marks on the rear of the mummy and tooth punctures in the skin indicate that Blue Babe was killed by a cave lion. Blue Babe appears to have died during the fall or winter, when it was relatively cold. The carcass probably cooled rapidly and soon froze, which made it difficult for scavengers to eat. Blue Babe is also frequently referenced when talking about scientists eating their own specimens: the research team that was preparing it for permanent display in the University of Alaska Museum removed a portion of the mummy's neck, stewed it, and dined on it to celebrate the accomplishment.
In early September 2007, near Tsiigehtchic, local resident Shane Van Loon discovered a carcass of a steppe bison which was radiocarbon dated to c. 13,650 cal BP. This carcass appears to represent the first Pleistocene mummified soft tissue remains from the glaciated regions of northern Canada.
In 2011, a 9,300-year-old mummy was found at Yukagir in Siberia.
In 2016, a frozen tail was discovered in the north of the Republic of Sakha in Russia. The exact age was not clear, but tests showed it was not younger than 8,000 years old. A team of Russian and South Korean scientists proposed extracting DNA from the specimen and cloning it in the future.
The steppe wisent is known from Denisova Cave, famous for being the site where the first Denisovan remains were discovered.
| Biology and health sciences | Bovidae | Animals |
21491359 | https://en.wikipedia.org/wiki/Quaoar | Quaoar | Quaoar (minor-planet designation: 50000 Quaoar) is a large, ringed dwarf planet in the Kuiper belt, a region of icy planetesimals beyond Neptune. It has an elongated ellipsoidal shape with an average diameter of , about half the size of the dwarf planet Pluto. The object was discovered by American astronomers Chad Trujillo and Michael Brown at the Palomar Observatory on 4 June 2002. Quaoar's surface contains crystalline water ice and ammonia hydrate, which suggests that it might have experienced cryovolcanism. A small amount of methane is present on its surface, which can only be retained by the largest Kuiper belt objects.
Quaoar has one known moon, Weywot, which was discovered by Brown in February 2007. Both objects were named after mythological figures from the Native American Tongva people in Southern California. Quaoar is the Tongva creator deity and Weywot is his son. In 2023, astronomers announced the discovery of two thin rings orbiting Quaoar outside its Roche limit, which defies theoretical expectations that rings outside the Roche limit should not be stable.
History
Discovery
Quaoar was discovered on 4 June 2002 by American astronomers Chad Trujillo and Michael Brown at the Palomar Observatory in the Palomar Mountain Range in San Diego County, California. The discovery formed part of the Caltech Wide Area Sky Survey, which was designed to search for the brightest Kuiper belt objects using the Palomar Observatory's 1.22-meter Samuel Oschin telescope. Quaoar was first identified in images by Trujillo on 5 June 2002, when he noticed a dim, 18.6-magnitude object slowly moving among the stars of the constellation Ophiuchus. Quaoar appeared relatively bright for a distant object, suggesting that it could have a size comparable to the diameter of Pluto.
To ascertain Quaoar's orbit, Brown and Trujillo initiated a search for archival precovery images. They obtained several precovery images taken by the Near-Earth Asteroid Tracking survey from various observatories in 1996 and 2000–2002. In particular, they had also found two archival photographic plates taken by astronomer Charles T. Kowal in May 1983, who at the time was searching for the hypothesized Planet X at the Palomar Observatory. From these precovery images, Brown and Trujillo were able to calculate Quaoar's orbit and distance. Additional precovery images of Quaoar have been later identified, with the earliest known found by Edward Rhoads on a photographic plate imaged on 25 May 1954 from the Palomar Observatory Sky Survey.
Before announcing the discovery of Quaoar, Brown had planned to conduct follow-up observations using the Hubble Space Telescope to measure Quaoar's size. He had also planned to announce the discovery as soon as possible and found it necessary to keep the discovery information confidential during the follow-up observations. Rather than submitting his Hubble proposal under peer review, Brown submitted his proposal directly to one of Hubble's operators, who promptly allocated time to Brown. While setting up the observing algorithm for Hubble, Brown had also planned to use one of the Keck telescopes in Mauna Kea, Hawaii, as a part of a study on cryovolcanism on the moons of Uranus. This provided him additional time for follow-up observations and took advantage of the whole observing session in July to analyze Quaoar's spectrum and characterize its surface composition.
The discovery of Quaoar was formally announced by the Minor Planet Center in a Minor Planet Electronic Circular on 7 October 2002. It was given the provisional designation , indicating that its discovery took place during the first half of June 2002. Quaoar was the 1,512th object discovered in the first half of June, as indicated by the preceding letter and numbers in its provisional designation. On that same day, Trujillo and Brown reported their scientific results from observations of Quaoar at the 34th annual meeting of the American Astronomical Society's Division for Planetary Sciences in Birmingham, Alabama. They announced Quaoar was the largest Kuiper belt object found yet, surpassing previous record holders 20000 Varuna and . Quaoar's discovery has been cited by Brown as having contributed to the reclassification of Pluto as a dwarf planet. Since then, Brown has contributed to the discovery of larger trans-Neptunian objects, including Haumea, , Makemake and .
Name and symbol
Upon Quaoar's discovery, it was initially given the temporary nickname "Object X" as a reference to Planet X, due to its potentially large size and unknown nature. At the time, Quaoar's size was uncertain, and its high brightness led the discovery team to speculate that it may be a possible tenth planet. After measuring Quaoar's size with the Hubble Space Telescope in July, the team began considering names for the object, particularly those from local Native American mythologies. Following the International Astronomical Union's (IAU) naming convention for minor planets, non-resonant Kuiper belt objects are to be named after creation deities. The team settled on the name Kwawar, the creator god of the Tongva people indigenous to the Los Angeles Basin, where Brown's institute, the California Institute of Technology, was located.
According to Brown, the name "Quaoar" is pronounced with three syllables, and Trujillo's website on Quaoar gives a three-syllable pronunciation, , as an approximation of the Tongva pronunciation . The name can be also pronounced as two syllables, , reflecting the usual English spelling and pronunciation of the deity Kwawar.
In Tongva mythology, Kwawar is the genderless creation force of the universe, singing and dancing deities into existence. He first sings and dances to create Weywot (Sky Father), then they together sing Chehooit (Earth Mother) and Tamit (Grandfather Sun) into existence. As they did this, the creation force became more complex as each new deity joined the singing and dancing. Eventually, after reducing chaos to order, they created the seven great giants that upheld the world, then the animals and finally the first man and woman, Tobohar and Pahavit.
Upon their investigation of names from Tongva mythology, Brown and Trujillo realized that there were contemporary members of the Tongva people, whom they contacted for permission to use the name. They consulted tribal historian Marc Acuña, who confirmed that the name Kwawar would indeed be an appropriate name for the newly discovered object. However, the Tongva preferred the spelling Qua-o-ar, which Brown and Trujillo adopted, though with the hyphens omitted. The name and discovery of Quaoar were publicly announced in October, though Brown had not sought approval of the name by the IAU's Committee on Small Body Nomenclature (CSBN). Indeed, Quaoar's name was announced before the official numbering of the object, which Brian Marsden—the head of the Minor Planet Center—remarked in 2004 to be a violation of the protocol. Despite this, the name was approved by the CSBN, and the naming citation, along with Quaoar's official numbering, was published in a Minor Planet Circular on 20 November 2002.
Quaoar was given the minor planet number 50000, which was not by coincidence but to commemorate its large size, being that it was found in the search for a Pluto-sized object in the Kuiper belt. The large Kuiper belt object 20000 Varuna was similarly numbered for a similar occasion. However, subsequent even larger discoveries such as 136199 Eris were simply numbered according to the order in which their orbits were confirmed.
The usage of planetary symbols is no longer recommended in astronomy, so Quaoar never received a symbol in the astronomical literature. A symbol , used mostly among astrologers, is included in Unicode as U+1F77E. The symbol was designed by Denis Moskowitz, a software engineer in Massachusetts; it combines the letter Q (for 'Quaoar') with a canoe, and is stylized to recall angular Tongva rock art.
Orbit and classification
Quaoar orbits the Sun at an average distance of , taking 288.8 years to complete one full orbit around the Sun. With an orbital eccentricity of 0.04, Quaoar follows a nearly circular orbit, only slightly varying in distance from 42 AU at perihelion to 45 AU at aphelion. At such distances, light from the Sun takes more than 5 hours to reach Quaoar. Quaoar has last passed aphelion in late 1932 and is currently approaching the Sun at a rate of 0.035 AU per year, or about . Quaoar will reach perihelion around February 2075.
Because Quaoar has a nearly circular orbit, it does not approach close to Neptune such that its orbit can become significantly perturbed under the gravitational influence of Neptune. Quaoar's minimum orbit intersection distance from Neptune is only 12.3 AU—it does not approach Neptune within this distance over the course of its orbit, as it is not in a mean-motion orbital resonance with Neptune. Simulations by the Deep Ecliptic Survey show that the perihelion and aphelion distances of Quaoar's orbit do not change significantly over the next ten million years; Quaoar's orbit appears to be stable over the long term.
Quaoar is a trans-Neptunian object. It is classified as a distant minor planet by the Minor Planet Center. Because Quaoar is not in a mean-motion resonance with Neptune, it is also classified as a classical Kuiper belt object (cubewano) by the Minor Planet Center and Deep Ecliptic Survey. Quaoar's orbit is moderately inclined to the ecliptic plane by 8 degrees, relatively high when compared to the inclinations of Kuiper belt objects within the dynamically cold population. Because Quaoar's orbital inclination is greater than 4 degrees, it is part of the dynamically hot population of high-inclination classical Kuiper belt objects. The high inclinations of hot classical Kuiper belt objects such as Quaoar are thought to have resulted from gravitational scattering by Neptune during its outward migration in the early Solar System.
Physical characteristics
Size and shape
, measurements of Quaoar's shape from its rotational light curve and stellar occultations show that Quaoar is a triaxial ellipsoid with an average diameter of . Quaoar's diameter is roughly half that of Pluto and is slightly smaller than Pluto's moon Charon. At the time of its discovery in 2002, Quaoar was the largest object found in the Solar System since the discovery of Pluto. Quaoar was also the first trans-Neptunian object to be measured directly from Hubble Space Telescope images.
Quaoar's far-infrared thermal emission and brightness in visible light both vary significantly (visible light curve amplitude 0.12–0.16 magnitudes) as Quaoar rotates every 17.68 hours, which most likely indicates Quaoar is elongated along its equator. A 2024 analysis of Quaoar's visible and far-infrared rotational light curve by Csaba Kiss and collaborators determined that the lengths of Quaoar's equatorial axes differ by 19% (a/b = 1.19) and the lengths of Quaoar's polar and shortest equatorial axis differ by 16% (b/c = 1.16), which corresponds to ellipsoid dimensions of . The ellipsoidal shape of Quaoar matches the size and shape measurements from previous stellar occultations, and also explains why the size and shape of Quaoar appeared to change in these occultations.
Quaoar's elongated shape contradicts theoretical expectations that it should be in hydrostatic equilibrium, because of its large size and slow rotation. According to Michael Brown, rocky bodies around in diameter should relax into hydrostatic equilibrium, whereas icy bodies relax into hydrostatic equilibrium somewhere between and . Slowly-rotating objects in hydrostatic equilibrium are expected to be oblate spheroids (Maclaurin spheroids), whereas rapidly-rotating objects in hydrostatic equilibrium, such as Haumea which rotates in nearly 4 hours, are expected to be flattened and elongated ellipsoids (Jacobi ellipsoids). To explain Quaoar's non-equilibrium shape, Kiss and collaborators hypothesized that Quaoar originally had a rapid rotation and was in hydrostatic equilibrium, but its shape became "frozen in" and did not change as Quaoar spun down due to tidal forces from its moon Weywot. This would resemble the situation of Saturn's moon Iapetus, which is too oblate for its current rotation rate.
Mass and density
Quaoar has a mass of , which was determined from Weywot's orbit using Kepler's third law. Measurements of Quaoar's diameter and mass indicate it has a density between , which suggests its interior is composed of roughly 70% rock and 30% ice with low porosity. Quaoar's density was previously thought to be much higher, between , because early measurements inaccurately suggested that Quaoar had a smaller diameter and a higher mass. These early high-density estimates for Quaoar led researchers to hypothesize that the object might be a rocky planetary core exposed by a large impact event, but these hypotheses have since become obsolete as newer estimates indicate a lower density for Quaoar.
Surface
Quaoar has a dark surface that reflects about 12% of the visible light it receives from the Sun. This may indicate that fresh ice has disappeared from Quaoar's surface. The surface is moderately red, meaning that Quaoar reflects longer (redder) wavelengths of light more than shorter (bluer) wavelengths. Many Kuiper belt objects such as 20000 Varuna and 28978 Ixion share a similar moderately red color.
Spectroscopic observations by David Jewitt and Jane Luu in 2004 revealed signs of crystalline water ice and ammonia hydrate on Quaoar's surface. These substances are expected to gradually break down due to solar and cosmic radiation, and crystalline water ice can only form in warm temperatures of at least , so the presence of crystalline water ice on Quaoar's surface indicates that it was heated to this temperature sometime in the last ten million years. For context, Quaoar's present-day surface temperature is less than . Jewitt and Luu proposed two hypotheses for Quaoar's heating, which are impact events and radiogenic heating. The latter hypothesis allows for the possibility of cryovolcanism on Quaoar, which is supported by the presence of ammonia hydrate on Quaoar's surface. Ammonia hydrate is believed to be cryovolcanically deposited onto Quaoar's surface. A 2006 study by Hauke Hussmann and collaborators suggested that radiogenic heating alone may not be capable of sustaining an internal ocean of liquid water at Quaoar's mantle–core boundary.
More precise observations of Quaoar's near infrared spectrum in 2007 indicated the presence of small quantities (5%) of solid methane and ethane. Given its boiling point of , methane is a volatile ice at average surface temperatures of Quaoar, unlike water ice or ethane. Both models and observations suggest that only a few larger bodies (Pluto, and ) can retain the volatile ices whereas the dominant population of small trans-Neptunian objects lost them. Quaoar, with only small amounts of methane, appears to be in an intermediary category.
In 2022, low-resolution near-infrared (0.7–5 μm) spectroscopic observations by the James Webb Space Telescope (JWST) revealed the presence of carbon dioxide ice, complex organics, and significant amounts of ethane ice on Quaoar's surface. Other possible chemical compounds include hydrogen cyanide and carbon monoxide. JWST also took medium-resolution near-infrared spectra of Quaoar and found evidence of small amounts of methane on Quaoar's surface. However, both JWST's low- and medium-resolution spectra of Quaoar did not show conclusive signs of ammonia hydrates.
Possible atmosphere
The presence of methane and other volatiles on Quaoar's surface suggest that it may support a tenuous atmosphere produced from the sublimation of volatiles. With a measured mean temperature of approximately , the upper limit of Quaoar's atmospheric pressure is expected to be in the range of a few microbars. Due to Quaoar's small size and mass, the possibility of Quaoar having an atmosphere of nitrogen and carbon monoxide has been ruled out, since the gases would escape from Quaoar. The possibility of a methane atmosphere, with the upper limit being less than 1 microbar, was considered until 2013, when Quaoar occulted a 15.8-magnitude star and revealed no sign of a substantial atmosphere, placing an upper limit to at least 20 nanobars, under the assumption that Quaoar's mean temperature is and that its atmosphere consists of mostly methane. The upper limit of atmosphere pressure was tightened to 10 nanobars after another stellar occultation in 2019.
Satellite
Quaoar has one known moon, Weywot (full designation (50000) Quaoar I Weywot), discovered in 2006 and named after the sky god Weywot, son of Quaoar. It orbits Quaoar at distance of about 13,300 km and is thought to be approximately in diameter.
Rings
Discovery
Besides accurately determining sizes and shapes, stellar occultation campaigns were planned on a long-term basis to search for rings and/or atmospheres around small bodies of the outer solar system. These campaigns agglomerated efforts of various teams in France, Spain and Brazil and were conducted under the umbrella of the European Research Council project Lucky Star. The discovery of Quaoar's first known ring, Q1R, involved various instruments used during stellar occultations observed between 2018 and 2021: the robotic ATOM telescope of the High Energy Stereoscopic System (HESS) in Namibia, the 10.4-m Gran Telescopio Canarias (La Palma Island, Spain); the ESA CHEOPS space telescope, and several stations run by citizen astronomers in Australia where a report of a Neptune-like ring originated and a dense arc in Q1R was first observed. Taken together, these observations reveal the presence of a partly dense, mostly tenuous and uniquely distant ring around Quaoar, a discovery announced in February 2023.
In April 2023, astronomers of the Lucky Star project published the discovery of another ring of Quaoar, Q2R. The Q2R ring was detected by the highly-sensitive 8.2-m Gemini North and the 4.0-m Canada-France-Hawaii Telescope in Mauna Kea, Hawaii, during an observing campaign to confirm Quaoar's Q1R ring in a stellar occultation on 9 August 2022. Quaoar is the fourth minor planet known and confirmed to have a ring system, after 10199 Chariklo, 2060 Chiron, and Haumea.
Properties
Quaoar possesses two narrow rings, provisionally named Q1R and Q2R by order of discovery, which are confined at radial distances where their orbital periods are integer ratios of Quaoar's rotational period. That is, the rings of Quaoar are in spin-orbit resonances.
The outer ring, Q1R, orbits Quaoar at a distance of , over seven times the radius of Quaoar and more than double the theoretical maximum distance of the Roche limit. The Q1R ring is not uniform and is strongly irregular around its circumference, being more opaque (and denser) where it is narrow and less opaque where it is broader. The Q1R ring's radial width ranges from while its optical depth ranges from 0.004 to 0.7. The irregular width of the Q1R ring resembles Saturn's frequently-perturbed F ring or Neptune's ring arcs, which may imply the presence of small, kilometer-sized moonlets embedded within the Q1R ring and gravitationally perturbing the material. The Q1R ring likely consists of icy particles that elastically collide with each other without accreting into a larger mass.
Q1R is located in between the 6:1 mean-motion orbital resonance with Quaoar's moon Weywot at and Quaoar's 1:3 spin-orbit resonance at . The Q1R ring's coincidental location at these resonances implies they play a key role in maintaining the ring without having it accrete into a single moon. In particular, the confinement of rings to the 1:3 spin-orbit resonance may be common among ringed small Solar System bodies, as it has been previously seen in Chariklo and Haumea.
The inner ring, Q2R, orbits Quaoar at a distance of , about four and a half times Quaoar's radius and also outside Quaoar's Roche limit. The Q2R ring's location coincides with Quaoar's 5:7 spin-orbit resonance at . Compared to Q1R, the Q2R ring appears relatively uniform with a radial width of . With an optical depth of 0.004, the Q2R ring is very tenuous and its opacity is comparable to the least dense part of the Q1R ring.
Exploration
It has been calculated that a flyby mission to Quaoar using a Jupiter gravity assist would take 13.6 years, for launch dates of 25 December 2026, 22 November 2027, 22 December 2028, 22 January 2030 and 20 December 2040. Quaoar would be 41 to 43 AU from the Sun when the spacecraft arrived. In July 2016, the Long Range Reconnaissance Imager (LORRI) aboard the New Horizons spacecraft took a sequence of four images of Quaoar from a distance of about 14 AU. Interstellar Probe, a concept by Pontus Brandt and his colleagues at Johns Hopkins Applied Physics Laboratory would potentially fly by Quaoar in the 2030s before continuing to the interstellar medium, and the first of China National Space Administration's proposed Shensuo probe designed to explore the heliosphere has it considered as a potential flyby target. Quaoar has been chosen as a flyby target for missions like these particularly for its escaping methane atmosphere and possible cryovolcanism, as well as its close proximity to the heliospheric nose.
| Physical sciences | Solar System | Astronomy |
21491745 | https://en.wikipedia.org/wiki/Argentine%20ant | Argentine ant | The Argentine ant (Linepithema humile, formerly Iridomyrmex humilis) is an ant native to northern Argentina, Uruguay, Paraguay, Bolivia and southern Brazil. This invasive species was inadvertently introduced by humans on a global scale and has become established in many Mediterranean climate areas, including South Africa, New Zealand, Japan, Easter Island, Australia, the Azores, Europe, Hawaii, and the continental United States. Argentine ants are significant pests within agricultural and urban settings, and are documented to cause substantial harm to communities of native arthropods, vertebrates, and plants within their invaded range.
Description
Linepithema humile is a small-bodied (2.2–2.6 mm) ant species, dull light to dark brown in color. Within the invasion zone, ant colonies are large and include many workers and multiple queens.
Argentine ants are opportunistic with regard to nesting preferences. Colony nests have been found in the ground, in cracks in concrete walls, in spaces between boards and timbers, even among belongings in human dwellings. In natural areas, they generally nest shallowly in loose leaf litter or beneath small stones, due to their poor ability to dig deeper nests. However, if a deeper nesting ant species abandons their nest, Argentine ant colonies will readily take over the space. Because the native habitat for this species is within riparian floodplains, colonies are very sensitive to water infiltration within their nests; if their nests become inundated with water, workers will collect the brood and the entire colony will move to dry ground.
Austrian entomologist Gustav L. Mayr identified the first specimens of Hypoclinea humilis in the vicinity of Buenos Aires, Argentina in 1866. This species was shortly transferred to the genus Iridomyrmex, and finally to Linepithema in the early 1990s.
Distribution
The native range of Argentine ants is limited to riparian habitats in the lowland areas of the Paraná River drainage, which stretches across northern Argentina, Uruguay, Paraguay, and southern Brazil. Within South America, this species has spread into parts of Chile, Colombia, Ecuador, and Peru. Linepithema humile thrives in Mediterranean climates, and over the past century it has spread to across the globe by human-mediated transport. The species has become established to every continent except Antarctica and includes many oceanic islands.
Global "mega-colony"
The absence of aggression within Argentine ant colonies was first reported in 1913 by Newell & Barber, who noted "…there is no apparent antagonism between separate colonies of its own kind".
Later studies showed that these "supercolonies" extend across hundreds or thousands of kilometers in different parts of the introduced range, first reported in California in 2000,
then in Europe in 2002,
Japan in 2009,
and Australia in 2010.
Several subsequent studies used genetic, behavioral, and chemical analyses to show that introduced supercolonies on separate continents actually represent a single global supercolony.
The researchers stated that the "enormous extent of this population is paralleled only by human society", and had probably been spread and maintained by human travel.
Behavior
They have been extraordinarily successful, in part, because different nests of the introduced Argentine ants seldom attack or compete with each other, unlike most other species of ant. In their introduced range, their genetic makeup is so uniform that individuals from one nest can mingle in a neighboring nest without being attacked. Thus, in most of their introduced range, they form supercolonies.
The Very Large Colony, which covers territory from San Diego to beyond San Francisco, may have a population of nearly one trillion individuals.
Conflict does occur between members of different supercolonies. In 1997, UC San Diego researchers observed fighting between different Argentine ants kept in lab, and in 2004 scientists began to map out the boundaries of the different supercolonies that clashed in San Diego. On the border of the Very Large Colony and the Lake Hodges Colony thirty million ants die each year, on a battlefront that covers many miles. While the battles of other ant species generally constitute colony raids lasting a few hours, or skirmishes that occur periodically for a few weeks, Argentine ants clash ceaselessly; the borders of their territory are a site of constant violence and battles can be fought on top of hundreds of dead ants. Fights may be halted by adverse weather such as rain.
In contrast, native populations are genetically more diverse and form colonies that are much smaller than the supercolonies that dominate the introduced range. Colonies living in close proximity are territorial and aggressive toward one another. Argentine ants in their native South America also co-exist with many other species of ants, and do not attain the high population densities that characterize introduced populations.
In a series of experiments, ants of the same colony were isolated and fed different diets. The hydrocarbons from the diet were eventually incorporated into the cuticle of the subjects. Those that had the same diet appeared to recognize one another as kin. Those who had at least some overlap in dietary composition also appeared to react non-aggressively to one another. These interactions contrasts drastically with the groups that fed on completely different sources, such as those who lived off flies and those that fed on grasshoppers. The groups appeared to have incorporated hydrocarbons that were not similar to the others and created an unfamiliar identity cue. These groups reacted violently towards each other. This suggests that dietary factors affect the recognition cues for colony members.
Reproduction and seasonal colony trends
Like workers in many other ant species, Argentine ant workers are unable to lay reproductive eggs but can direct the development of eggs into reproductive females; the production of males appears to be controlled by the amount of food available to the larvae. Argentine ant colonies almost invariably have many reproductive queens, as many as eight for every 1,000 workers.
The seasonal low occurs in mid-winter, when 90% of a representative colony consists of workers and the remainder of queens, and no reproductive activity and minimal birthing. Eggs are produced in late-winter, nearly all of which hatch into sexual forms by May(*). Mating occurs after the females emerge. Worker production increases steadily from mid-March(*) to October(*), after which their numbers are not replenished; thus, their numbers drop steadily over the winter months.
((*) Note that the information regarding the months May, March, and October in this paragraph, as well as the entire Month axis in the graph, are most likely not worldwide correct, in special in Southern hemisphere, due to the shift of six months in seasons between the two hemispheres.)
Colonies in the Argentine ant's native habitat are kept within a range of ten to one hundred meters by colonies of interspecific and intraspecific rivals. As the colonies expand, they appear to form fluctuating territory borders, which contract and expand on a seasonal and conditional basis. There is an expansive push outward in the summer months, with a retreating motion in the winter. This has to do with soil moisture and temperature conditions.
At the edges of these borders are either rival L. humile colonies or other obstacles that prevent further expansion, such as an inhospitable environment for nests.
Impact
The ants are ranked among the world's 100 worst invasive animal species.
In its introduced range, the Argentine ant often displaces most or all native ants and can threaten native invertebrates and even small vertebrates that are not accustomed to defending against the aggressive ants. This can, in turn, imperil other species in the ecosystem, such as native plants that depend on native ants for seed dispersal, or lizards that depend on native ants or invertebrates for food. For example, the recent severe decline in coastal horned lizards in southern California is closely tied to Argentine ants displacing native ant species on which the lizards feed. In South Africa, the Argentine ant has in some cases displaced native ants that disperse the seeds of Fynbos plants like Mimetes cucullatus. The Argentine ants don't take the seeds underground and are left on the surface, resulting in ungerminated plants and the dwindling of Fynbos seed reserves after veld fires.
Argentine ants sometimes tend aphid, mealybug, and scale insect colonies,
sometimes relocating the parasites to unaffected plants, and their protection of these plant pests from predators and parasitoids can cause problems in agricultural areas.
In return for this protection, the ants benefit by feeding off an excretion known as "honeydew".
There is also evidence that the presence of Argentine ant may decrease the number of pollinators that visit natural flowering plants via predation on the larvae of the pollinators.
Pest control
Argentine ants are a common household pest, often entering structures in search of food or water (particularly during dry or hot weather), or to escape flooded nests during periods of heavy rainfall. When they invade a kitchen, it is not uncommon to see two or three queens foraging along with the workers.
Borate-sucrose water baits are toxic to Argentine ants, when the bait is 25% water, with 0.5–1.0% boric acid or borate salts.
In spring, during a colony's growth phase, protein based baits may be more effective due to much higher demand from the egg-laying queens.
Due to their nesting behavior and presence of numerous queens in each colony, it is generally impractical to spray Argentine ants with pesticides or to use boiling water as with mound building ants. Spraying with pesticides has occasionally stimulated increased egg-laying by the queens, compounding the problem. Pest control usually requires exploiting their omnivorous dietary habits, through use of slow-acting poison bait (e.g. fipronil, hydramethylnon, sulfluramid), which will be carried back to the nest by the workers, eventually killing all the individuals, including the queens.
Research
Researchers from the University of California, Irvine, have developed a way to use the scent of Argentine ants against them. The exoskeletons of the ants are covered with a hydrocarbon-laced secretion. They made a compound that is different, but similar, to the one that coats the ants. If the chemical is applied to an ant, the other members of the colony will kill it.
Another approach for a large scale control of the Argentine ant has been proposed by researchers from Japan, who showed that it is possible to disrupt its trails with synthetic pheromones. This has been confirmed in various later trials by a New Zealand-led team in Hawaii and by researchers from Victoria University of Wellington who showed that this approach is beneficial for other local ant species.
A multi-agency study at the environmentally-sensitive California Channel Islands in the mid-2010s eliminated the Argentine ant to sub-detectable levels. The ant was introduced to the islands in the early- to mid-1900s. Due to strict biosecurity measures, the risk of re-invasion was low and allowed for a controlled study. Four (4) km2, 2% of the island area, was treated with liquid bait pellets similar to off-the-shelf products, dispersed by helicopter. Cost of treatment ranged from $50000-$200000 USD (2015) per km2 including the cost of subsequent monitoring, and may serve as an initial guideline for other areas.
| Biology and health sciences | Hymenoptera | Animals |
21492663 | https://en.wikipedia.org/wiki/Asbestos | Asbestos | Asbestos ( ) is a group of naturally occurring, toxic, carcinogenic and fibrous silicate minerals. There are six types, all of which are composed of long and thin fibrous crystals, each fibre (particulate with length substantially greater than width) being composed of many microscopic "fibrils" that can be released into the atmosphere by abrasion and other processes. Inhalation of asbestos fibres can lead to various dangerous lung conditions, including mesothelioma, asbestosis, and lung cancer. As a result of these health effects, asbestos is considered a serious health and safety hazard.
Archaeological studies have found evidence of asbestos being used as far back as the Stone Age to strengthen ceramic pots, but large-scale mining began at the end of the 19th century when manufacturers and builders began using asbestos for its desirable physical properties. Asbestos is an excellent thermal and electrical insulator, and is highly fire resistant, so for much of the 20th century, it was very commonly used around the world as a building material (particularly for its fire-retardant properties), until its adverse effects on human health were more widely recognized and acknowledged in the 1970s. Many buildings constructed before the 1980s contain asbestos.
The use of asbestos for construction and fireproofing has been made illegal in many countries. Despite this, around 255,000 people are thought to die each year from diseases related to asbestos exposure. In part, this is because many older buildings still contain asbestos; in addition, the consequences of exposure can take decades to arise. The latency period (from exposure until the diagnosis of negative health effects) is typically 20 years. The most common diseases associated with chronic asbestos exposure are asbestosis (scarring of the lungs due to asbestos inhalation) and mesothelioma (a type of cancer).
Many developing countries still support the use of asbestos as a building material, and mining of asbestos is ongoing, with the top producer, Russia, having an estimated production of 790,000 tonnes in 2020.
Etymology
The word "asbestos", first used in the 1600s, ultimately derives from the , meaning "unquenchable" or "inextinguishable". The name reflects use of the substance for wicks that would never burn up.
It was adopted into English via the Old French abestos, which got the word from Greek via Latin, but in the original Greek, "asbestos" actually referred to quicklime. It is said by the Oxford English Dictionary that the word was wrongly used by Pliny for what we now call asbestos, and that he popularized the misnomer. Asbestos was originally referred to in Greek as amiantos, meaning "undefiled", because when thrown into a fire it came out unmarked. "Amiantos" is the source for the word for asbestos in many languages, such as the Portuguese amianto and the French amiante. It had also been called "amiant" in English in the early 15th century, but this usage was superseded by "asbestos". The word is pronounced or .
History
Asbestos has been used for thousands of years to create flexible objects that resist fire, including napkins, but, in the modern era, companies began producing consumer goods containing asbestos on an industrial scale. Today, the risk of asbestos has been recognized; the use of asbestos is completely banned in 66 countries and strictly regulated in many others.
Early references and uses
Asbestos use dates back at least 4,500 years, when the inhabitants of the Lake Juojärvi region in East Finland strengthened earthenware pots and cooking utensils with the asbestos mineral anthophyllite; archaeologists call this style of pottery "asbestos-ceramic". Some archaeologists believe that ancient peoples made shrouds of asbestos, wherein they burned the bodies of their kings to preserve only their ashes and to prevent the ashes being mixed with those of wood or other combustible materials commonly used in funeral pyres. Others assert that these peoples used asbestos to make perpetual wicks for sepulchral or other lamps. A famous example is the golden lamp asbestos lychnis, which the sculptor Callimachus made for the Erechtheion. In more recent centuries, asbestos was indeed used for this purpose.
A once-purported first description of asbestos occurs in Theophrastus, On Stones, from around 300 BC, but this identification has been refuted. In both modern and ancient Greek, the usual name for the material known in English as "asbestos" is amiantos ("undefiled", "pure"), which was adapted into the French as amiante and into Italian, Spanish and Portuguese as amianto. In modern Greek, the word ἀσβεστος or ασβέστης stands consistently and solely for lime.
The term asbestos is traceable to Roman naturalist Pliny the Elder's first-century manuscript Natural History and his use of the term asbestinon, meaning "unquenchable"; he described the mineral as being more expensive than pearls. While Pliny or his nephew Pliny the Younger is popularly credited with recognising the detrimental effects of asbestos on human beings, examination of the primary sources reveals no support for either claim.
In China, accounts of obtaining huo huan bu () or "fire-laundered cloth" certainly dates to the Wei dynasty, and there are claims they were known as early as in the Zhou dynasty, though not substantiated except in writings thought to date much later (see huoshu or "fire rat").
Athanasius of Alexandria, a Christian bishop living in 4th-century Egypt, references asbestos in one of his writings. Around the year 318, he wrote as follows:
Wealthy Persians amazed guests by cleaning a cloth by exposing it to fire. For example, according to Tabari, one of the curious items belonging to Khosrow II Parviz, the great Sassanian king (r. 590–628), was a napkin () that he cleaned simply by throwing it into fire. Such cloth is believed to have been made of asbestos imported over the Hindu Kush. According to Biruni in his book Gems, any cloths made of asbestos (, āzarshost) were called shostakeh (). Some Persians believed the fiber was the fur of an animal called the samandar (), which lived in fire and died when exposed to water; this was where the former belief originated that the salamander could tolerate fire. Charlemagne, the first Holy Roman Emperor (800–814), is also said to have possessed such a tablecloth.
Marco Polo recounts having been shown, in a place he calls Ghinghin talas, "a good vein from which the cloth which we call of salamander, which cannot be burnt if it is thrown into the fire, is made ..."
Industrial era
The large-scale asbestos industry began in the mid-19th century. Early attempts at producing asbestos paper and cloth in Italy began in the 1850s but were unsuccessful in creating a market for such products. Canadian samples of asbestos were displayed in London in 1862, and the first companies were formed in England and Scotland to exploit this resource. Asbestos was first used in the manufacture of yarn, and German industrialist Louis Wertheim adopted this process in his factories in Germany.
In 1871, the Patent Asbestos Manufacturing Company was established in Glasgow, and during the following decades, the Clydebank area became a centre for the nascent industry.
Industrial-scale mining began in the Thetford hills, Quebec, from the 1870s. Sir William Edmond Logan was the first to notice the large deposits of chrysotile in the hills in his capacity as head of Geological Survey of Canada. Samples of the minerals from there were displayed in London and elicited much interest. With the opening of the Quebec Central Railway in 1876, mining entrepreneurs such as Andrew Stuart Johnson established the asbestos industry in the province. The 50-ton output of the mines in 1878 rose to over 10,000 tonnes in the 1890s with the adoption of machine technologies and expanded production. For a long time, the world's largest asbestos mine was the Jeffrey mine in the town of Asbestos, Quebec.
Asbestos production began in the Urals of the Russian Empire in the 1880s, and the Alpine regions of Northern Italy with the formation in Turin of the Italo-English Pure Asbestos Company in 1876, although this was soon swamped by the greater production levels from the Canadian mines. Mining also took off in South Africa from 1893 under the aegis of the British businessman Francis Oates, the director of the De Beers company. It was in South Africa that the production of amosite began in 1910. The U.S. asbestos industry had an early start in 1858 when fibrous anthophyllite was mined for use as asbestos insulation by the Johns Company, a predecessor to the current Johns Manville, at a quarry at Ward's Hill on Staten Island, New York. US production began in earnest in 1899 with the discovery of large deposits in Belvidere Mountain in Vermont.
The use of asbestos became increasingly widespread toward the end of the 19th century when its diverse applications included fire-retardant coatings, concrete, bricks, pipes and fireplace cement, heat-, fire-, and acid-resistant gaskets, pipe insulation, ceiling insulation, fireproof drywall, flooring, roofing, lawn furniture, and drywall joint compound. In 2011, it was reported that over 50% of UK houses still contained asbestos, despite a ban on asbestos products some years earlier.
In Japan, particularly after World War II, asbestos was used in the manufacture of ammonium sulfate for purposes of rice production, sprayed upon the ceilings, iron skeletons, and walls of railroad cars and buildings (during the 1960s), and used for energy efficiency reasons as well. Production of asbestos in Japan peaked in 1974 and went through ups and downs until about 1990 when production began to drop dramatically.
Discovery of toxicity
The 1898 Annual Report of the Chief Inspector of Factories and Workshops in the United Kingdom noted the negative health effects of asbestos, the contribution having been made by Lucy Deane Streatfeild, one of the first women factory inspectors.
In 1899, H. Montague Murray noted the negative health effects of asbestos. The first documented death related to asbestos was in 1906.
In the early 1900s, researchers began to notice a large number of early deaths and lung problems in asbestos-mining towns. The first such study was conducted by Murray at the Charing Cross Hospital, London, in 1900, in which a postmortem investigation discovered asbestos traces in the lungs of a young man who had died from pulmonary fibrosis after having worked for 14 years in an asbestos textile factory. Adelaide Anderson, the Inspector of Factories in Britain, included asbestos in a list of harmful industrial substances in 1902. Similar investigations were conducted in France in 1906 and Italy in 1908.
The first diagnosis of asbestosis was made in the UK in 1924. Nellie Kershaw was employed at Turner Brothers Asbestos in Rochdale, Greater Manchester, England, from 1917, spinning raw asbestos fibre into yarn. Her death in 1924 led to a formal inquest. Pathologist William Edmund Cooke testified that his examination of the lungs indicated old scarring indicative of a previous, healed tuberculosis infection, and extensive fibrosis, in which were visible "particles of mineral matter ... of various shapes, but the large majority have sharp angles." Having compared these particles with samples of asbestos dust provided by S. A. Henry, His Majesty's Medical Inspector of Factories, Cooke concluded that they "originated from asbestos and were, beyond a reasonable doubt, the primary cause of the fibrosis of the lungs and therefore of death."
As a result of Cooke's paper, Parliament commissioned an inquiry into the effects of asbestos dust by E. R. A. Merewether, Medical Inspector of Factories, and , a factory inspector and pioneer of dust monitoring and control. Their subsequent report, Occurrence of Pulmonary Fibrosis & Other Pulmonary Affections in Asbestos Workers, was presented to Parliament on 24 March 1930. It concluded that the development of asbestosis was irrefutably linked to the prolonged inhalation of asbestos dust, and included the first health study of asbestos workers, which found that 66% of those employed for 20 years or more suffered from asbestosis. The report led to the publication of the first asbestos industry regulations in 1931, which came into effect on 1 March 1932. These rules regulated ventilation and made asbestosis an excusable work-related disease. The term mesothelioma was first used in medical literature in 1931; its association with asbestos was first noted sometime in the 1940s. Similar legislation followed in the U.S. about ten years later.
Approximately 100,000 people in the United States have died, or are terminally ill, from asbestos exposure related to shipbuilding. In the Hampton Roads area, a shipbuilding center, mesothelioma occurrence is seven times the national rate. Thousands of tons of asbestos were used in World War II ships to insulate piping, boilers, steam engines, and steam turbines. There were approximately 4.3 million shipyard workers in the United States during the war; for every 1,000 workers, about 14 died of mesothelioma and an unknown number died of asbestosis.
The United States government and the asbestos industry have been criticized for not acting quickly enough to inform the public of dangers and to reduce public exposure. In the late 1970s, court documents proved that asbestos industry officials knew of asbestos dangers since the 1930s and had concealed them from the public.
In Australia, asbestos was widely used in construction and other industries between 1946 and 1980. From the 1970s, there was increasing concern about the dangers of asbestos, and following community and union campaigning, its use was phased out, with mining having ceased in 1983. The use of asbestos was phased out in 1989 and banned entirely in December 2003. The dangers of asbestos are now well known in Australia, and there is help and support for those suffering from asbestosis or mesothelioma.
Use by industry and product type
Serpentine group
Serpentine minerals have a sheet or layered structure. Chrysotile (commonly known as white asbestos) is the only asbestos mineral in the serpentine group. In the United States, chrysotile has been the most commonly used type of asbestos. According to the U.S. Environmental Protection Agency (EPA) Asbestos Building Inspectors Manual, chrysotile accounts for approximately 95% of asbestos found in buildings in the United States. Chrysotile is often present in a wide variety of products and materials, including:
Chlor Alkali diaphragm membranes used to make chlorine (currently in the US)
Drywall and joint compound (including texture coats)
Plaster
Gas mask filters throughout World War II until the 1960s for most countries; Germany and the USSR's Civilian issued filters up until 1988 tested positive for asbestos
Vinyl floor tiles, sheeting, adhesives
Roofing tars, felts, siding, and shingles
"Transite" panels, siding, countertops, and pipes
Popcorn ceilings, also known as acoustic ceilings
Fireproofing
Caulk
Industrial and marine gaskets
Brake pads and shoes
Stage curtains
Fire blankets
Cement pipework
Interior fire doors
Fireproof clothing for firefighters
Thermal pipe insulation
Filters for removing fine particulates from chemicals, liquids, and wine
Dental cast linings
HVAC flexible duct connectors
Drilling fluid additives
In the European Union and Australia, it has been banned as a potential health hazard and is no longer used at all.
Amphibole group
Amphiboles including amosite (brown asbestos) and crocidolite (blue asbestos) were formerly used in many products until the early 1980s. Tremolite asbestos constituted a contaminant of many if not all naturally occurring chrysotile deposits. The use of all types of asbestos in the amphibole group was banned in much of the Western world by the mid-1980s, and in Japan by 1995. Some products that included amphibole types of asbestos included the following:
Low-density insulating board (often referred to as AIB or asbestos insulating board) and ceiling tiles;
Asbestos cement sheets and pipes for construction, casing for water and electrical/telecommunication services;
Thermal and chemical insulation (e.g., fire-rated doors, limpet spray, lagging, and gaskets).
Electrical wiring, braided cables, cable wrap, wire insulation (usually crocidolite)
Cigarette manufacturer Lorillard (Kent's filtered cigarette) used crocidolite asbestos in its "Micronite" filter from 1952 to 1956.
While mostly chrysotile asbestos fibers were once used in automobile brake pads, shoes, and clutch discs, contaminants of amphiboles were present. Since approximately the mid-1990s, brake pads, new or replacement, have been manufactured instead with linings made of ceramic, carbon, metallic, and aramid fiber (Twaron or Kevlar—the same material used in bulletproof vests).
Artificial Christmas snow, known as flocking, was previously made with asbestos. It was used as an effect in films including The Wizard of Oz and department store window displays and it was marketed for use in private homes under brand names that included "Pure White", "Snow Drift" and "White Magic".
Potential use in carbon sequestration
The potential for use of asbestos to mitigate climate change has been raised. Although the adverse aspects of mining of minerals, including health effects, must be taken into account, exploration of the use of mineral wastes to sequester carbon is being studied. The use of mining waste materials from nickel, copper, diamond, and platinum mines have the potential as well, but asbestos may have the greatest potential and is the subject of research now in progress in an emerging field of scientific study to examine it. The most common type of asbestos, chrysotile, chemically reacts with CO2 to produce ecologically stable magnesium carbonate. Chrysotile, like all types of asbestos, has a large surface area that provides more places for chemical reactions to occur, compared to most other naturally occurring materials.
Construction
Developed countries
The use of asbestos in new construction projects has been banned for health and safety reasons in many developed countries or regions, including the European Union, the United Kingdom, Australia, Hong Kong, Japan, and New Zealand. A notable exception is the United States, where asbestos continues to be used in construction such as cement asbestos pipes. The 5th Circuit Court prevented the EPA from banning asbestos in 1991 because EPA research showed the ban would cost between US$450 and 800 million while only saving around 200 lives in a 13-year timeframe, and that the EPA did not provide adequate evidence for the safety of alternative products. Until the mid-1980s, small amounts of white asbestos were used in the manufacture of Artex, a decorative stipple finish, however, some of the lesser-known suppliers of Artex-type materials were still adding white asbestos until 1999.
Before the ban, asbestos was widely used in the construction industry in thousands of materials. Some are judged to be more dangerous than others due to the amount of asbestos and the material's friable nature. Sprayed coatings, pipe insulation, and Asbestos Insulating Board (AIB) are thought to be the most dangerous due to their high content of asbestos and friable nature. Many older buildings built before the late 1990s contain asbestos. In the United States, there is a minimum standard for asbestos surveys as described by ASTM standard E 2356–18. In the UK, the Health and Safety Executive have issued guidance called HSG264 describing how surveys should be completed although other methods can be used if they can demonstrate they have met the regulations by other means. The EPA includes some, but not all, asbestos-contaminated facilities on the Superfund National Priorities List (NPL). Renovation and demolition of asbestos-contaminated buildings are subject to EPA NESHAP and OSHA Regulations. Asbestos is not a material covered under CERCLA's innocent purchaser defense. In the UK, the removal and disposal of asbestos and substances containing it are covered by the Control of Asbestos Regulations 2006.
U.S. asbestos consumption hit a peak of 804,000 tons in 1973; world asbestos demand peaked around 1977, with 25 countries producing nearly 4.8 million metric tons annually.
In older buildings (e.g. those built before 1999 in the UK, before white asbestos was banned), asbestos may still be present in some areas. Being aware of asbestos locations reduces the risk of disturbing asbestos.
Removal of asbestos building components can also remove the fire protection they provide, therefore fire protection substitutes are required for proper fire protection that the asbestos originally provided.
Outside Europe and North America
Some countries, such as India, Indonesia, China and Russia, have continued widespread use of asbestos. The most common is corrugated asbestos-cement sheets or "A/C sheets" for roofing and sidewalls. Millions of homes, factories, schools or sheds, and shelters continue to use asbestos. Cutting these sheets to size and drilling holes to receive 'J' bolts to help secure the sheets to roof framing is done on-site. There has been no significant change in production and use of A/C sheets in developing countries following the widespread restrictions in developed nations.
September 11 attacks
As New York City's World Trade Center collapsed following the September 11 attacks, Lower Manhattan was blanketed in a mixture of building debris and combustible materials. This complex mixture gave rise to the concern that thousands of residents and workers in the area would be exposed to known hazards in the air and dust, such as asbestos, lead, glass fibers, and pulverized concrete. More than 1,000 tons of asbestos are thought to have been released into the air following the buildings' destruction. Inhalation of a mixture of asbestos and other toxicants is thought to be linked to the unusually high death rate from cancer of emergency service workers since the disaster. Thousands more are now thought to be at risk of developing cancer due to this exposure with those who have died so far being only the "tip of the iceberg".
In May 2002, after numerous cleanup, dust collection, and air monitoring activities were conducted outdoors by EPA, other federal agencies, New York City, and the state of New York, New York City formally requested federal assistance to clean and test residences in the vicinity of the World Trade Center site for airborne asbestos.
Asbestos contaminants in other products
Vermiculite
Vermiculite is a hydrated laminar magnesium-aluminum-iron silicate that resembles mica. It can be used for many industrial applications and has been used as insulation. Some deposits of vermiculite are contaminated with small amounts of asbestos.
One vermiculite mine operated by W. R. Grace and Company in Libby, Montana exposed workers and community residents to danger by mining vermiculite contaminated with asbestos, typically richterite, winchite, actinolite or tremolite. Vermiculite contaminated with asbestos from the Libby mine was used as insulation in residential and commercial buildings through Canada and the United States. W. R. Grace and Company's loose-fill vermiculite was marketed as Zonolite but was also used in sprayed-on products such as Monokote.
In 1999, the EPA began cleanup efforts in Libby and now the area is a Superfund cleanup area. The EPA has determined that harmful asbestos is released from the mine as well as through other activities that disturb soil in the area.
Talc
Talc can sometimes be contaminated with asbestos due to the proximity of asbestos ore (usually tremolite) in underground talc deposits. By 1973, US federal law required all talc products to be asbestos-free. Separating cosmetic-grade talc (e.g. talcum powder) from industrial-grade talc (often used in friction products) has largely eliminated this issue for consumers. Cosmetics companies, including Johnson & Johnson, have known since the 1950s that talc products could be contaminated with asbestos. In 2020, laboratory tests of 21 talc-based cosmetics products found that 15 percent were contaminated with asbestos. In July 2024, the World Health Organization (WHO) heightened its health warning about talc exposure. In a study published in The Lancet, the WHO changed its classification of talc from "possibly carcinogenic" to "probably carcinogenic."
In 2000, tests in a certified asbestos-testing laboratory found the tremolite form of amphibole asbestos used to be found in three out of eight popular brands of children's crayons that were made partly from talc: Crayola, Prang, and RoseArt. In Crayola crayons, the tests found asbestos levels around 0.05% in Carnation Pink and 2.86% in Orchid; in Prang crayons, the range was from 0.3% in Periwinkle to 0.54% in Yellow; in Rose Art crayons, it was from 0.03% in Brown to 1.20% in Orange. Overall, 32 different types of crayons from these brands used to contain more than trace amounts of asbestos, and eight others contained trace amounts. The Art and Creative Materials Institute, a trade association which tested the safety of crayons on behalf of the makers, initially insisted the test results must have been incorrect, although they later said they do not test for asbestos. In May 2000, Crayola said tests by Richard Lee, a materials analyst whose testimony on behalf of the asbestos industry has been accepted in lawsuits over 250 times, found its crayons tested negative for asbestos. In spite of that, in June 2000 Binney & Smith, the maker of Crayola, and the other makers agreed to stop using talc in their products, and changed their product formulations in the United States.
The mining company R T Vanderbilt Co of Gouverneur, New York, which supplied the talc to the crayon makers, states that "to the best of our knowledge and belief" there had never been any asbestos-related disease among the company's workers. However media reports claim that the United States Mine Safety and Health Administration (MSHA) had found asbestos in four talc samples tested in 2000. The Assistant Secretary for Mine Safety and Health subsequently wrote to the news reporter, stating that "In fact, the abbreviation ND (non-detect) in the laboratory report – indicates no asbestos fibers actually were found in the samples." Multiple studies by mineral chemists, cell biologists, and toxicologists between 1970 and 2000 found neither samples of asbestos in talc products nor symptoms of asbestos exposure among workers dealing with talc, but more recent work has rejected these conclusions in favor of "same as" asbestos risk.
On 12 July 2018, a Missouri jury ordered Johnson & Johnson to pay a record $4.69 billion to 22 women who alleged the company's talc-based products, including its baby powder, contain asbestos and caused them to develop ovarian cancer.
Types and associated fibers
Six mineral types are defined by the EPA as "asbestos" including those belonging to the serpentine class and those belonging to the amphibole class. All six asbestos mineral types are known to be human carcinogens. The visible fibers are themselves each composed of millions of microscopic "fibrils" that can be released by abrasion and other processes.
Serpentine class: chrysotile
Serpentine class fibers are curly. Chrysotile, CAS No. , is the only asbestos classed as a serpentine fiber. It is obtained from serpentinite rocks which are common throughout the world. Its idealized chemical formula is Mg(SiO)(OH). Chrysotile appears under the microscope as a white fiber.
Chrysotile has been used more than any other type and accounts for about 95% of the asbestos found in buildings in America. Chrysotile is more flexible than amphibole types of asbestos and can be spun and woven into fabric. The most common use was corrugated asbestos cement roofing primarily for outbuildings, warehouses, and garages. It may also be found in sheets or panels used for ceilings and sometimes for walls and floors. Chrysotile has been a component in joint compound and some plasters. Numerous other items have been made containing chrysotile including brake linings, fire barriers in fuseboxes, pipe insulation, floor tiles, residential shingles, and gaskets for high-temperature equipment.
Amphibole class
Amphibole class fibers are needle-like. Amosite, crocidolite, tremolite, anthophyllite and actinolite are members of the amphibole class.
Amosite
Amosite, CAS No. , often referred to as brown asbestos, is a trade name for the amphiboles belonging to the cummingtonite-grunerite solid solution series, commonly from South Africa, named as a partial acronym for "Asbestos Mines of South Africa". One formula given for amosite is . Amosite is seen under a microscope as a grey-white vitreous fiber. It is found most frequently as a fire retardant in thermal insulation products, asbestos insulating board and ceiling tiles.
Crocidolite
Crocidolite, CAS No. , commonly known as blue asbestos, is the fibrous form of the amphibole riebeckite, found primarily in southern Africa, but also in Australia and Bolivia. One formula given for crocidolite is . Crocidolite is seen under a microscope as a blue fiber.
Crocidolite commonly occurs as soft friable fibers. Asbestiform amphibole may also occur as soft friable fibers but some varieties such as amosite are commonly straighter. All forms of asbestos are fibrillar in that they are composed of fibers with breadths less than 1 micrometer in bundles of very great widths. Asbestos with particularly fine fibers is also referred to as "amianthus".
Tremolite, anthophyllite, actinolite
Other regulated asbestos minerals, such as tremolite asbestos, CAS No. , ; actinolite asbestos, CAS No. 77536-66-4, ; and anthophyllite asbestos, CAS No. , ; are less commonly used industrially but can still be found in a variety of construction materials and insulation materials and have been used in a few consumer products.
Other natural asbestiform minerals, such as richterite, , and winchite, , though not regulated, are said by some to be no less harmful than tremolite, amosite, or crocidolite. They are termed "asbestiform" rather than asbestos. Although the U.S. Occupational Safety and Health Administration (OSHA) has not included them in the asbestos standard, NIOSH and the American Thoracic Society have recommended them for inclusion as regulated materials because they may also be hazardous to health.
"Mountain leather" is an old-fashioned term for flexible, sheet-like natural formations of asbestiform minerals which resemble leather. Asbestos-containing minerals known to form mountain leather include: actinolite, sepiolite, and tremolite.
Production
In 2017, 1.3 million tonnes of asbestos were mined worldwide. Russia was the largest producer with 53% of the world total, followed by Kazakhstan (16%), China (15%), and Brazil (11.5%). Asia consumes some 70% of the asbestos produced in the world with China, India and Indonesia the largest consumers.
In 2009, about 9% of the world's asbestos production was mined in Canada. In late 2011, Canada's remaining two asbestos mines, both located in Quebec, halted operations. In September 2012, the Quebec government halted asbestos mining.
Health impact
The most common diseases associated with chronic asbestos exposure are asbestosis (scarring of the lungs due to asbestos inhalation) and mesothelioma (cancer associated with asbestos). Mesothelioma is an aggressive form of cancer and often leads to a life expectancy of less than 12 months after diagnosis.
All types of asbestos fibers are known to represent serious health hazards in humans and animals. Amosite and crocidolite are considered the most hazardous asbestos fiber types; however, chrysotile asbestos has also produced tumors in animals and is a recognized cause of asbestosis and malignant mesothelioma in humans, and mesothelioma has been observed in people who were occupationally exposed to chrysotile, family members of the occupationally exposed, and residents who lived close to asbestos factories and mines.
During the 1980s and again in the 1990s, the asbestos industry suggested at times that the process of making asbestos cement could "neutralize" the asbestos, either via chemical processes or by causing the cement to attach to the fibers and changing their physical size; subsequent studies showed that this was untrue and that decades-old asbestos cement, when broken, releases asbestos fibers identical to those found in nature, with no detectable alteration.
Exposure to asbestos in the form of fibers is always considered dangerous. Working with, or exposure to, material that is friable, or materials or works that could cause the release of loose asbestos fibers, is considered high risk. In general, people who become ill from inhaling asbestos have been regularly exposed in a job where they worked directly with the material.
The US Occupational Safety and Health Administration (OSHA) has standards to protect workers from the hazards of exposure to asbestos in the workplace. The permissible exposure limit for asbestos is 0.1 fiber per cubic centimeter of air as an eight-hour time-weighted average, with an excursion limit of 1.0 asbestos fibers per cubic centimeter over a 30-minute period.
Regulation
Complete bans on asbestos
Worldwide, 66 countries and territories (including all those in the European Union) have banned the use of asbestos. Exemptions for minor uses are permitted in some countries listed; however, all countries listed must have banned the use of all types of asbestos.
Australia
The use of crocidolite (blue asbestos) was banned in 1967, while the use of amosite (brown asbestos) continued in the construction industry until the mid-1980s. It was finally banned from building products in 1989, though it remained in gaskets and brake linings until 31 December 2003, and cannot be imported, used, or recycled.
Asbestos continues to be a problem in Australia. Two out of three homes in Australia built between World War II and the early 1980s still contain asbestos.
The United Services Union, representing workers tasked with modifying electrical meter boxes at residences, stated that workers should refuse to do this work until the boxes have been inspected for asbestos, and the head of the Australian Council of Trade Unions (ACTU) has called on the government to protect its citizens by ridding the country of asbestos by 2030.
Handlers of asbestos materials must have a B-Class license for bonded asbestos and an A-Class license for friable asbestos.
The town of Wittenoom, in Western Australia, was built around a (blue) asbestos mine. The entire town continues to be contaminated and has been disincorporated, allowing local authorities to remove references to Wittenoom from maps and road signs.
In January 2024, asbestos was found in garden mulch supplied to dozens of sites including parks, playgrounds and schools across Sydney, triggering the Sydney asbestos mulch crisis.
Canada
As of December 31, 2018, it is illegal to import, manufacture, sell, trade, or use products made from asbestos. There are exemptions for its use in the chloralkali industry, the military, nuclear facilities, and for magnesium extraction from asbestos mining residues.
Iran
On 23 July 2000 Supreme Council of Environmental Protection of Iran banned any use of asbestos. On a note it states that after 44 years if there is no replacement technically, economically or financially the law is appealable and will be subject to review.
On 22 November 2011 Iran's Environmental Protection Organization banned use of asbestos completely, it also banned all imports and exports of asbestos. On 9 October 2011 Supreme Council of Environmental Protection of Iran confirmed the new law and publicly announced it, updating and replacing the previous law.
Japan
Revelations that hundreds of workers had died in Japan over the previous few decades from diseases related to asbestos sparked a scandal in mid-2005. Tokyo had, in 1971, ordered companies handling asbestos to install ventilators and check health regularly; however, the Japanese government did not ban crocidolite and amosite until 1995, and a near-complete ban with a few exceptions on asbestos was implemented in 2006, with the remaining exceptions being removed in March 2012 for a full-fledged ban.
New Zealand
In 1984, the import of raw amphibole (blue and brown) asbestos into New Zealand was banned. In 2002 the import of chrysotile (white) asbestos was also banned. In 2015, the government announced that the importation of asbestos would be completely banned with very limited exceptions (expected to be applied to replacement parts for older machines) that would be reviewed on a case-by-case basis.
North-west of Nelson, in the Upper Takaka Valley, is New Zealand's only commercially harvested asbestos mine. A low-grade Chrysotile was mined here from 1908 to 1917 but only 100 tons were washed and taken out by packhorse. A new power scheme enabled work to renew and between 1940 and 1949, 40 tons a month was mined by the Hume Company. This continued to 1964, when, due to the short length of its fibre, the limited commercial viability forced mining to cease.
South Korea
In May 1997, the manufacture and use of crocidolite and amosite, commonly known as blue and brown asbestos, were fully banned in South Korea. In January 2009, a full-fledged ban on all types of asbestos occurred when the government banned the manufacture, import, sale, storage, transport or use of asbestos or any substance containing more than 0.1% of asbestos. In 2011, South Korea became the world's sixth country to enact an asbestos harm aid act, which entitles any Korean citizen to free lifetime medical care as well as monthly income from the government if they are diagnosed with an asbestos-related disease.
United Kingdom
In the United Kingdom, blue and brown asbestos materials were banned outright in 1985 while the import, sale, and secondhand reuse of white asbestos was outlawed in 1999. The 2012 Control of Asbestos Regulations, updating and replacing the previous 2006 law, state that owners of non-domestic buildings (e.g., factories and offices) have a "duty to manage" asbestos on the premises by making themselves aware of its presence and ensuring the material does not deteriorate, removing it if necessary. Employers, e.g. construction companies, whose operatives may come into contact with asbestos must also provide annual asbestos training to their workers.
United States
Prior to 2024, the United States remained one of the few developed countries to not completely ban asbestos. Some American workers at chlorine plants frequently came in contact with the substance, and OSHA exempted these plants from random inspections through the Voluntary Protection Program.
In 1989, the United States Environmental Protection Agency (EPA) issued the Asbestos Ban and Phase-Out Rule, but in 1991, asbestos industry supporters challenged and overturned the ban in a landmark lawsuit: Corrosion Proof Fittings v. the Environmental Protection Agency. Although the case resulted in several small victories for asbestos regulation, the EPA ultimately did not put an end to asbestos use. The ruling left many consumer products that can still legally contain trace amounts of asbestos. Six categories of asbestos-containing products are however banned: corrugated paper, rollboard, commercial paper, specialty paper, flooring felt and any new uses of asbestos. The Clean Air Act also bans asbestos pipe insulation and asbestos block insulation on components such as boilers and hot water tanks, and spray-applied surfacing asbestos-containing materials. The Consumer Product Safety Act bans asbestos in artificial fireplace embers and wall patching compounds. The Food and Drug Administration bans asbestos-containing filters in pharmaceutical manufacturing, processing, and packing.
In 2014, the state of Washington banned asbestos in automotive brakes.
In 2022, 224 metric tons of chrysotile asbestos were imported into the United States from Brazil and Russia, for use in diaphragms used to produce chlorine by the chloralkali process. As of 2020, imported goods containing asbestos included gaskets for some utility vehicles, rubber sheets used in gasket manufacturing, brake blocks for the oil industry, and vehicle friction materials such as aftermarket automotive brakes and linings.
In 2024, the EPA announced a new rule to ban all ongoing uses of asbestos by 2037. The rule is the first to be implemented under 2016 amendments to the Toxic Substances Control Act.
Countries where asbestos is legal
Mexico
Since 1970, as a result of increased regulation of asbestos in Europe and in the United States, there was a massive transfer of asbestos-processing enterprises to Mexico. Asbestos is used in many products – roofing, boilers, pipes, brakes, and wires, produced by over 2,000 Mexican companies, many of them subsidiaries or subcontractors of US companies, and sold throughout the Americas. In 2000, 58% of Mexican asbestos-containing exports went to the United States, and 40% to Central American countries and Cuba.
Vietnam
In Vietnam, chrysotile asbestos is not banned and is still widely used. Amphibole asbestos is banned from trade and use. Vietnam is one of the top 10 asbestos users in the world, with an annual import volume of about 65,000–70,000 tons of chrysotile. About 90% of the imported asbestos is used to produce about 100 million m2 of cement roofing sheets (asbestos-cement). According to one study, among 300 families in Yen Bai, Thanh Hoa, 85% of households use asbestos roofing sheets, but only 5% know about the negative health effects.
The master plan (for construction materials development to 2020 with orientation to 2030 submitted by the Ministry of Construction to the Government in January 2014) still suggests continued use of chrysotile for a long time.
Substitutes for asbestos in construction
Fiberglass insulation was invented in 1938 and is now the most commonly used type of insulation material. The safety of this material has also been called into question due to similarities in material structure. However, the International Agency for Research on Cancer removed fiberglass from its list of possible human carcinogens in 2001. A scientific review article from 2011 claimed epidemiology data was inconsistent and concluded that the IARC's decision to downgrade the carcinogenic potential of fiberglass was valid, although this study was funded by a sponsored research contract from the North American Insulation Manufacturer's Association.
In 1978, a highly texturized fiberglass fabric was invented by Bal Dixit, called Zetex. This fabric is lighter than asbestos but offers the same bulk, thickness, hand, feel, and abrasion resistance as asbestos. The fiberglass was texturized to eliminate some of the problems that arise with fiberglass, such as poor abrasion resistance and poor seam strength.
In Europe, mineral wool and glass wool are the main insulators in houses.
Many companies that produced asbestos-cement products that were reinforced with asbestos fibers have developed products incorporating organic fibers. One such product was known as "Eternit" and another "Everite" now use "Nutec" fibers which consist of organic fibers, portland cement and silica. Cement-bonded wood fiber is another substitute. Stone fibers are used in gaskets and friction materials.
Another potential fiber is polybenzimidazole or PBI fiber. Polybenzimidazole fiber is a synthetic fiber with a high melting point of that also does not ignite. Because of its exceptional thermal and chemical stability, it is often used by fire departments and space agencies.
Recycling and disposal
In most developed countries, asbestos is typically disposed of as hazardous waste in designated landfill sites.
The demolition of buildings containing large amounts of asbestos-based materials pose particular problems for builders and property developers – such buildings often have to be deconstructed piece by piece, or the asbestos has to be painstakingly removed before the structure can be razed by mechanical or explosive means. One such example is the Red Road Flats in Glasgow, Scotland which used huge amounts of asbestos cement board for wall panelling – British health and safety regulations stipulate that asbestos material has to be removed in specially adapted vehicles and taken to a landfill site with an appropriate permit to accept asbestos, via an approved route, at certain times of the day.
In the United States, the EPA governs the removal and disposal of asbestos strictly. Companies that remove asbestos must comply with EPA licensing. These companies are called EPA licensed asbestos contractors. Anytime one of these asbestos contractors performs work a test consultant has to conduct strict testing to ensure the asbestos is completely removed.
Asbestos can be destroyed by ultra-high-temperature incineration and plasma melting process. A process of thermal decomposition at produces a mixture of non-hazardous silicon-based wastes, and at temperatures above it produces silicate glass. Microwave thermal treatment can be used in an industrial manufacturing process to transform asbestos and asbestos-containing waste into porcelain stoneware tiles, porous single-fired wall tiles, and ceramic bricks.
The combination of oxalic acid with ultrasound fully degrades chrysotile asbestos fibers.
Abbreviations associated with asbestos
ACM: Asbestos-containing material (technically, material containing more than 1% asbestos)
AIB: Asbestos insulating board
| Physical sciences | Mineralogy | null |
21492980 | https://en.wikipedia.org/wiki/IMac | IMac | The iMac is a series of all-in-one computers from Apple Inc., sold as part of the company's Mac family of computers. First introduced in 1998, it has remained a primary part of Apple's consumer desktop offerings since and evolved through seven distinct forms. The iMac natively runs the MacOS operating system.
In its original form, the iMac G3 had a gumdrop, ADM-3 or egg-shaped look, with a CRT monitor, mainly enclosed by a colored, translucent plastic case. The computer was, at the time, an inexpensive, consumer-oriented computer that would easily connect to the Internet. The second major revision, the iMac G4, moved a design with a hemispherical base containing all the main components and an LCD monitor on a freely moving arm attached to it. The third and fourth revisions, the iMac G5 and the Intel iMac, placed all the components immediately behind the display in a plastic casing, creating a slim unified design that tilts only up and down on a simple metal base. The fifth, sixth and seventh revisions swapped the plastic enclosure for metal and became progressively thinner over each revision.
The design of the iMac has been seen as both controversial and trendsetting. From its introduction, the computer has eschewed many entrenched legacy technologies, notably becoming an early adopter of the USB port, and removing floppy disk and later optical disc drives. The most recent revision, the Apple Silicon iMac, uses Apple's own processors (silicon) and is thick. Between 2017 and 2021, Apple also sold a workstation-class version of the computer called the iMac Pro.
History
Apple was facing bankruptcy in the mid-1990s, with its market share cannibalized by Windows-based PCs and Macintosh clones. The company had tried and failed to ship a modern operating system for its hardware. Looking instead for an outside product to acquire, Apple announced its purchase of NexT, Inc. in 1996. Alongside Next's products and software came Steve Jobs, Apple's co-founder who had been ousted from the company years earlier. Jobs initially was brought on at Apple as an adviser, but Jobs replaced Gil Amelio as interim CEO in 1997 and began a reorganization of the company. He reduced Apple's multitude of confusing computer options to just four: one laptop and one desktop model for consumers, and another laptop and desktop model for professionals. What became the iMac began as Apple's effort to develop the consumer desktop to fill that product gap.
Apple's head of design Jony Ive and the rest of the design team developed sketches for a distinctive, all-in-one computer that was to be a legacy-free PC focused on ease of use and internet connectivity. The design team made the new computer colorful and translucent, built around a cathode-ray tube display wrapped in a curved plastic case. Ad agency director Ken Segall suggested the "iMac" name: it was short, had "Mac" in it, and the "i" prefix suggested the internet. Jobs initially hated it, but the name ultimately stuck. Apple later adopted the 'i' prefix across its consumer hardware and software lines, such as iPod, iBook (later MacBook), iPhone, iPad and various pieces of software such as the iLife, iCloud suite and iWork and the company's media player/store, iTunes.
Despite mixed reviews from the tech press, the iMac was a major commercial success at a time when Apple desperately needed a hit product. The iMac ultimately sold more than six million units, being revised multiple times and appearing in 13 different colors and patterns. The iMac was "designed to make it easy for home users to connect to the Internet." A commercial, dubbed "Simplicity Shootout", pitted seven-year-old Johann Thomas and his border collie Brodie, with an iMac, against Adam Taggart, a Stanford University MBA student, with an HP Pavilion 8250, in a race to set up their computers. Johann and Brodie finished in 8 minutes and 15 seconds, whereas Adam was still working on it by the end of the commercial.
As the prices of flat-screen liquid crystal displays (LCDs) began to fall, Apple conceived of an update to the iMac. Inspired by a sunflower, the iMac G4 put the computer in a semi-hemispherical base, with the display sitting above it on a stainless steel arm. The arm allowed the display to be easily tilted, rotated, and raised and lowered by a touch. The exuberant colors of the old iMac was replaced by stark white.
Ever-increasing screen sizes led Apple to make the iMac G5 a more conservative design, with the components of the computer attached to the back of the display and raised above the resting surface with an aluminum foot.
By 2005, it had become more and more apparent that IBM's development for the desktop implementation of PowerPC was grinding to a halt. Apple announced at the Worldwide Developers Conference that it would be switching the Macintosh to the x86 architecture and Intel's line of Core processors. The first Intel-equipped Macs were unveiled on January 10, 2006: the MacBook Pro and a new iMac, which outwardly looked identical to the iMac G5. Within nine months, Apple had smoothly transitioned the entire Macintosh line to Intel. The Intel-based iMac was redesigned in 2007 with an aluminum enclosure, which was gradually refined and slimmed down in the following years. In 2014, the iMac added high-resolution "retina" 4K and 5K displays, and a more powerful, professional-oriented model, the iMac Pro, was introduced in 2017.
Apple announced a shift from Intel processors to its own Apple Silicon in June 2020. Apple announced redesigned iMacs with a 24-inch display and Apple M1 chip in April 2021. These new models harkened back to the colorful iMac G3s, coming in seven colors. The iMacs were updated in 2023 to use the Apple M3 chip.
Influence
The original iMac was the first legacy-free PC. It was the first Macintosh computer to have a USB port but no floppy disk drive. Subsequently, all Macs have included USB. Via the USB port, hardware makers could make products compatible with both x86 PCs and Macs. Previously, Macintosh users had to seek out certain hardware, such as keyboards and mice specifically tailored for the "old world" Mac's unique ADB interface and printers and modems with MiniDIN-8 serial ports. Only a limited number of models from certain manufacturers were made with these interfaces and often came at a premium price. USB, being cross-platform, has allowed Macintosh users to select from a large selection of devices marketed for the Wintel PC platform, such as hubs, scanners, storage devices, USB flash drives, and mice. After the iMac, Apple continued to remove older peripheral interfaces and floppy drives from the rest of its product line.
Borrowing from the 1997 Twentieth Anniversary Macintosh, the various LCD-based iMac designs continued the all-in-one concept first envisioned in Apple's original Macintosh computer. The successful iMac allowed Apple to continue targeting the Power Macintosh line at the high-end of the market. This foreshadowed a similar strategy in the notebook market when the iMac-like iBook was released in 1999. Since then, the company has continued this strategy of differentiating the consumer versus professional product lines. Apple's focus on design has allowed each of its subsequent products to create a distinctive identity. Apple avoided using the beige colors that were then common in the PC industry. The company would later drift from the multicolored designs of the late 1990s and early 2000s. The latter part of the decade saw Apple using anodized aluminum; glass; and white, black, and clear polycarbonate plastics among its build materials. Today many PCs are more design-conscious than before the iMac's introduction, with multi-shaded design schemes being common, and some desktops and laptops available in colorful, decorative patterns.
Apple's use of translucent, candy-colored plastics inspired similar industrial designs in other consumer products. Apple's later introduction of the iPod, iBook G3 (Dual USB), and iMac G4 (all featuring snowy-white plastic), inspired similar designs in other companies' consumer electronics products. The color rollout also featured two distinctive ads: one called 'Life Savers' featured the Rolling Stones song, "She's a Rainbow" and an advertisement for the white version had the introduction of Cream's "White Room" as its backing track.
Reception
iMac has received considerable critical acclaim, including praise from technology columnist Walt Mossberg as the "Gold Standard of desktop computing"; Forbes magazine described the original candy-colored line of iMac computers as being an "industry-altering success". The first 24" Core 2 Duo iMac received CNET's "Must-have desktop" in its 2006 Top 10 Holiday Gift Picks.
Apple faced a class-action lawsuit filed in 2008 for allegedly deceiving the public by promising millions of colors from the LCD screens of all Mac models while its 20-inch model only held 262,144 colors. This issue arose due to the use of 6-bit per pixel Twisted nematic LCD screens. The case was dismissed on January 21, 2009.
While not a criticism of iMac per se, the integrated design has some inherent tradeoffs that have garnered criticism. In The Mythical Midrange Mac Minitower, Dan Frakes of Macworld suggests that with the iMac occupying the midrange of Apple's product line, Apple has little to offer consumers who want some ability to expand or upgrade their computers, but do not need (or cannot afford) the Mac Pro. For example, iMac's integration of monitor and computer, while convenient, commits the owner to replace both at the same time. For a time before the Mac mini's introduction, there were rumors of a "headless iMac" but the G4 Mac mini as introduced had lower performance compared to the iMac, which at the time featured a G5 processor. Some third party suppliers such as Other World Computing provide upgrade kits that include specialized tools for working on iMacs.
Similarly, though the graphics chipset in some Intel models is on a removable MXM, neither Apple nor third parties have offered retail iMac GPU upgrades, with the exception of those for the original iMac G3's "mezzanine" PCI slot. Models after iMac G5 (excluding the August 7, 2007, iMac update) made it difficult for the end-user to replace the hard disk or optical drive, and Apple's warranty explicitly forbids upgrading the socketed CPU. While conceding the possibility of a mini-tower cannibalizing sales from the Mac Pro, Frakes argues there is enough frustration with iMac's limitations to make such a proposition worthwhile. This disparity has become more pronounced after the G4 era since the bottom-end Power Mac G5 (with one brief exception) and Mac Pro models have all been priced in the US$1999–$2499 range, while base model Power Macs G4s and earlier were US$1299–1799. The current generation iMac has Intel 5th generation i5 and i7 processors, ranging from quad-core 2.7 GHz i5 to a quad-core 3.4 GHz i7 processor, however it is possible to upgrade the 2010 edition of the iMac quite easily.
Timeline
Supported operating systems
Supported Apple operating system releases
macOS Sequoia is the current release of macOS, being compatible with 2019 or later iMacs. Most unsupported Intel iMac computers can run macOS Sequoia via the use of a compatible utility.
Supported Windows versions
| Technology | Specific hardware | null |
3864143 | https://en.wikipedia.org/wiki/Igneous%20differentiation | Igneous differentiation | In geology, igneous differentiation, or magmatic differentiation, is an umbrella term for the various processes by which magmas undergo bulk chemical change during the partial melting process, cooling, emplacement, or eruption. The sequence of (usually increasingly silicic) magmas produced by igneous differentiation is known as a magma series.
Definitions
Primary melts
When a rock melts to form a liquid, the liquid is known as a primary melt. Primary melts have not undergone any differentiation and represent the starting composition of a magma. In nature, primary melts are rarely seen. Some leucosomes of migmatites are examples of primary melts. Primary melts derived from the mantle are especially important and are known as primitive melts or primitive magmas. By finding the primitive magma composition of a magma series, it is possible to model the composition of the rock from which a melt was formed, which is important because we have little direct evidence of the Earth's mantle.
Parental melts
Where it is impossible to find the primitive or primary magma composition, it is often useful to attempt to identify a parental melt. A parental melt is a magma composition from which the observed range of magma chemistries has been derived by the processes of igneous differentiation. It need not be a primitive melt.
For instance, a series of basalt lava flows is assumed to be related to one another. A composition from which they could reasonably be produced by fractional crystallization is termed a parental melt. To prove this, fractional crystallization models would be produced to test the hypothesis that they share a common parental melt.
Cumulate rocks
Fractional crystallization and accumulation of crystals formed during the differentiation process of a magmatic event are known as cumulate rocks, and those parts are the first which crystallize out of the magma. Identifying whether a rock is a cumulate or not is crucial for understanding if it can be modelled back to a primary melt or a primitive melt, and identifying whether the magma has dropped out cumulate minerals is equally important even for rocks which carry no phenocrysts.
Underlying causes of differentiation
The primary cause of change in the composition of a magma is cooling, which is an inevitable consequence of the magma being formed and migrating from the site of partial melting into an area of lower stress - generally a cooler volume of the crust.
Cooling causes the magma to begin to crystallize minerals from the melt or liquid portion of the magma. Most magmas are a mixture of liquid rock (melt) and crystalline minerals (phenocrysts).
Contamination is another cause of magma differentiation. Contamination can be caused by assimilation of wall rocks, mixing of two or more magmas or even by replenishment of the magma chamber with fresh, hot magma.
The whole gamut of mechanisms for differentiation has been referred to as the FARM process, which stands for fractional crystallization, assimilation, replenishment and magma mixing.
Fractional crystallization of igneous rocks
Fractional crystallization is the removal and segregation from a melt of mineral precipitates, which changes the composition of the melt. This is one of the most important geochemical and physical processes operating within the Earth's crust and mantle.
Fractional crystallization in silicate melts (magmas) is a very complex process compared to chemical systems in the laboratory because it is affected by a wide variety of phenomena. Prime amongst these are the composition, temperature, and pressure of a magma during its cooling.
The composition of a magma is the primary control on which mineral is crystallized as the melt cools down past the liquidus. For instance in mafic and ultramafic melts, the MgO and SiO2 contents determine whether forsterite olivine is precipitated or whether enstatite pyroxene is precipitated.
Two magmas of similar composition and temperature at different pressure may crystallize different minerals. An example is high-pressure and high-temperature fractional crystallization of granites to produce single-feldspar granite, and low-pressure low-temperature conditions which produce two-feldspar granites.
The partial pressure of volatile phases in silicate melts is also of prime importance, especially in near-solidus crystallization of granites.
Assimilation
Assimilation can be broadly defined as a process where a mass of magma wholly or partially homogenizes with materials derived from the wall rock of the magma body. Assimilation is a popular mechanism to partly explain the felsification of ultramafic and mafic magmas as they rise through the crust: a hot primitive melt intruding into a cooler, felsic crust will melt the crust and mix with the resulting melt. This then alters the composition of the primitive magma. Also, pre-existing mafic host rocks can be assimilated by very hot primitive magmas.
Effects of assimilation on the chemistry and evolution of magma bodies are to be expected, and have been clearly proven in many places. In the early 20th century there was a lively discussion on the relative importance of the process in igneous differentiation. More recent research has shown, however, that assimilation has a fundamental role in altering the trace element and isotopic composition of magmas, in formation of some economically important ore deposits, and in causing volcanic eruptions.
Replenishment
When a melt undergoes cooling along the liquid line of descent, the results are limited to the production of a homogeneous solid body of intrusive rock, with uniform mineralogy and composition, or a partially differentiated cumulate mass with layers, compositional zones and so on. This behaviour is fairly predictable and easy enough to prove with geochemical investigations. In such cases, a magma chamber will form a close approximation of the ideal Bowen's reaction series. However, most magmatic systems are polyphase events, with several pulses of magmatism. In such a case, the liquid line of descent is interrupted by the injection of a fresh batch of hot, undifferentiated magma. This can cause extreme fractional crystallisation because of three main effects:
Additional heat provides additional energy to allow more vigorous convection, allows resorption of existing mineral phases back into the melt, and can cause a higher-temperature form of a mineral or other higher-temperature minerals to begin precipitating
Fresh magma changes the composition of the melt, changing the chemistry of the phases which are being precipitated. For instance, plagioclase conforms to the liquid line of descent by forming initial anorthite which, if removed, changes the equilibrium mineral composition to oligoclase or albite. Replenishment of the magma can see this trend reversed, so that more anorthite is precipitated atop cumulate layers of albite.
Fresh magma destabilises minerals which are precipitating as solid solution series or on a eutectic; a change in composition and temperature can cause extremely rapid crystallisation of certain mineral phases which are undergoing a eutectic crystallisation phase.
Magma mixing
Magma mixing is the process by which two magmas meet, comingle, and form a magma of a composition somewhere between the two end-member magmas.
Magma mixing is a common process in volcanic magma chambers, which are open-system chambers where magmas enter the chamber, undergo some form of assimilation, fractional crystallisation and partial melt extraction (via eruption of lava), and are replenished.
Magma mixing also tends to occur at deeper levels in the crust and is considered one of the primary mechanisms for forming intermediate rocks such as monzonite and andesite. Here, due to heat transfer and increased volatile flux from subduction, the silicic crust melts to form a felsic magma (essentially granitic in composition). These granitic melts are known as an underplate. Basaltic primary melts formed in the mantle beneath the crust rise and mingle with the underplate magmas, the result being part-way between basalt and rhyolite; literally an 'intermediate' composition.
Other mechanisms of differentiation
Interface entrapment
Convection in a large magma chamber is subject to the interplay of forces generated by thermal convection and the resistance offered by friction, viscosity and drag on the magma offered by the walls of the magma chamber. Often near the margins of a magma chamber which is convecting, cooler and more viscous layers form concentrically from the outside in, defined by breaks in viscosity and temperature. This forms laminar flow, which separates several domains of the magma chamber which can begin to differentiate separately.
Flow banding is the result of a process of fractional crystallization which occurs by convection, if the crystals which are caught in the flow-banded margins are removed from the melt. The friction and viscosity of the magma causes phenocrysts and xenoliths within the magma or lava to slow down near the interface and become trapped in a viscous layer. This can change the composition of the melt in large intrusions, leading to differentiation.
Partial melt extraction
With reference to the definitions, above, a magma chamber will tend to cool down and crystallize minerals according to the liquid line of descent. When this occurs, especially in conjunction with zonation and crystal accumulation, and the melt portion is removed, this can change the composition of a magma chamber. In fact, this is basically fractional crystallization, except in this case we are observing a magma chamber which is the remnant left behind from which a daughter melt has been extracted.
If such a magma chamber continues to cool, the minerals it forms and its overall composition will not match a sample liquid line of descent or a parental magma composition.
Typical behaviours of magma chambers
It is worth reiterating that magma chambers are not usually static single entities. The typical magma chamber is formed from a series of injections of melt and magma, and most are also subject to some form of partial melt extraction.
Granite magmas are generally much more viscous than mafic magmas and are usually more homogeneous in composition. This is generally considered to be caused by the viscosity of the magma, which is orders of magnitude higher than mafic magmas. The higher viscosity means that, when melted, a granitic magma will tend to move in a larger concerted mass and be emplaced as a larger mass because it is less fluid and able to move. This is why granites tend to occur as large plutons, and mafic rocks as dikes and sills.
Granites are cooler and are therefore less able to melt and assimilate country rocks. Wholesale contamination is therefore minor and unusual, although mixing of granitic and basaltic melts is not unknown where basalt is injected into granitic magma chambers.
Mafic magmas are more liable to flow, and are therefore more likely to undergo periodic replenishment of a magma chamber. Because they are more fluid, crystal precipitation occurs much more rapidly, resulting in greater changes by fractional crystallisation. Higher temperatures also allow mafic magmas to assimilate wall rocks more readily and therefore contamination is more common and better developed.
Dissolved gases
All igneous magmas contain dissolved gases (water, carbonic acid, hydrogen sulfide, chlorine, fluorine, boric acid, etc.). Of these water is the principal, and was formerly believed to have percolated downwards from the Earth's surface to the heated rocks below, but is now generally admitted to be an integral part of the magma. Many peculiarities of the structure of the plutonic rocks as contrasted with the lavas may reasonably be accounted for by the operation of these gases, which were unable to escape as the deep-seated masses slowly cooled, while they were promptly given up by the superficial effusions. The acid plutonic or intrusive rocks have never been reproduced by laboratory experiments, and the only successful attempts to obtain their minerals artificially have been those in which special provision was made for the retention of the "mineralizing" gases in the crucibles or sealed tubes employed. These gases often do not enter into the composition of the rock-forming minerals, for most of these are free from water, carbonic acid, etc. Hence as crystallization goes on the residual melt must contain an ever-increasing proportion of volatile constituents. It is conceivable that in the final stages the still uncrystallized part of the magma has more resemblance to a solution of mineral matter in superheated steam than to a dry igneous fusion. Quartz, for example, is the last mineral to form in a granite. It bears much of the stamp of the quartz which we know has been deposited from aqueous solution in veins, etc. It is at the same time the most infusible of all the common minerals of rocks. Its late formation shows that in this case it arose at comparatively low temperatures and points clearly to the special importance of the gases of the magma as determining the sequence of crystallization.
When solidification is nearly complete the gases can no longer be retained in the rock and make their escape through fissures towards the surface. They are powerful agents in attacking the minerals of the rocks which they traverse, and instances of their operation are found in the kaolinization of granites, tourmalinization and formation of greisen, deposition of quartz veins, and the group of changes known as propylitization. These "pneumatolytic" processes are of the first importance in the genesis of many ore deposits. They are a real part of the history of the magma itself and constitute the terminal phases of the volcanic sequence.
Quantifying igneous differentiation
There are several methods of directly measuring and quantifying igneous differentiation processes;
Whole rock geochemistry of representative samples, to track changes and evolution of the magma systems
Using the above, calculating normative mineralogy and investigating trends
Trace element geochemistry
Isotope geochemistry
Investigating the contamination of magma systems by wall rock assimilation using radiogenic isotopes
In all cases, the primary and most valuable method for identifying magma differentiation processes is mapping the exposed rocks, tracking mineralogical changes within the igneous rocks and describing field relationships and textural evidence for magma differentiation. Clinopyroxene thermobarometry can be used to determine pressures and temperatures of magma differentiation.
| Physical sciences | Geochemistry | Earth science |
3871014 | https://en.wikipedia.org/wiki/Rainbow | Rainbow | A rainbow is an optical phenomenon caused by refraction, internal reflection and dispersion of light in water droplets resulting in a continuous spectrum of light appearing in the sky. The rainbow takes the form of a multicoloured circular arc. Rainbows caused by sunlight always appear in the section of sky directly opposite the Sun. Rainbows can be caused by many forms of airborne water. These include not only rain, but also mist, spray, and airborne dew.
Rainbows can be full circles. However, the observer normally sees only an arc formed by illuminated droplets above the ground, and centered on a line from the Sun to the observer's eye.
In a primary rainbow, the arc shows red on the outer part and violet on the inner side. This rainbow is caused by light being refracted when entering a droplet of water, then reflected inside on the back of the droplet and refracted again when leaving it.
In a double rainbow, a second arc is seen outside the primary arc, and has the order of its colours reversed, with red on the inner side of the arc. This is caused by the light being reflected twice on the inside of the droplet before leaving it.
Visibility
Rainbows can be observed whenever there are water drops in the air and sunlight shining from behind the observer at a low altitude angle. Because of this, rainbows are usually seen in the western sky during the morning and in the eastern sky during the early evening. The most spectacular rainbow displays happen when half the sky is still dark with raining clouds and the observer is at a spot with clear sky in the direction of the Sun. The result is a luminous rainbow that contrasts with the darkened background. During such good visibility conditions, the larger but fainter secondary rainbow is often visible. It appears about 10° outside of the primary rainbow, with inverse order of colours.
The rainbow effect is also commonly seen near waterfalls or fountains. In addition, the effect can be artificially created by dispersing water droplets into the air during a sunny day. Rarely, a moonbow, lunar rainbow or nighttime rainbow, can be seen on strongly moonlit nights. As human visual perception for colour is poor in low light, moonbows are often perceived to be white.
It is difficult to photograph the complete semicircle of a rainbow in one frame, as this would require an angle of view of 84°. For a 35 mm camera, a wide-angle lens with a focal length of 19 mm or less would be required. Now that software for stitching several images into a panorama is available, images of the entire arc and even secondary arcs can be created fairly easily from a series of overlapping frames.
From above the Earth such as in an aeroplane, it is sometimes possible to see a rainbow as a full circle. This phenomenon can be confused with the glory phenomenon, but a glory is usually much smaller, covering only 5–20°.
The sky inside a primary rainbow is brighter than the sky outside of the bow. This is because each raindrop is a sphere and it scatters light over an entire circular disc in the sky. The radius of the disc depends on the wavelength of light, with red light being scattered over a larger angle than blue light. Over most of the disc, scattered light at all wavelengths overlaps, resulting in white light which brightens the sky. At the edge, the wavelength dependence of the scattering gives rise to the rainbow.
The light of a primary rainbow arc is 96% polarised tangential to the arc. The light of the second arc is 90% polarised.
Number of colours in a spectrum or a rainbow
For colours seen by the human eye, the most commonly cited and remembered sequence is Isaac Newton's sevenfold red, orange, yellow, green, blue, indigo and violet, remembered by the mnemonic Richard Of York Gave Battle In Vain, or as the name of a fictional person (Roy G. Biv). The initialism is sometimes referred to in reverse order, as VIBGYOR. More modernly, the rainbow is often divided into red, orange, yellow, green, cyan, blue and violet. The apparent discreteness of main colours is an artefact of human perception and the exact number of main colours is a somewhat arbitrary choice.
Newton, who admitted his eyes were not very critical in distinguishing colours, originally (1672) divided the spectrum into five main colours: red, yellow, green, blue and violet. Later he included orange and indigo, giving seven main colours by analogy to the number of notes in a musical scale. Newton chose to divide the visible spectrum into seven colours out of a belief derived from the beliefs of the ancient Greek sophists, who thought there was a connection between the colours, the musical notes, the known objects in the Solar System, and the days of the week. Scholars have noted that what Newton regarded at the time as "blue" would today be regarded as cyan, and what Newton called "indigo" would today be considered blue.
The colour pattern of a rainbow is different from a spectrum, and the colours are less saturated. There is spectral smearing in a rainbow since, for any particular wavelength, there is a distribution of exit angles, rather than a single unvarying angle. In addition, a rainbow is a blurred version of the bow obtained from a point source, because the disk diameter of the sun (0.533°) cannot be neglected compared to the width of a rainbow (2.36°). Further red of the first supplementary rainbow overlaps the violet of the primary rainbow, so rather than the final colour being a variant of spectral violet, it is actually a purple. The number of colour bands of a rainbow may therefore be different from the number of bands in a spectrum, especially if the droplets are particularly large or small. Therefore, the number of colours of a rainbow is variable. If, however, the word rainbow is used inaccurately to mean spectrum, it is the number of main colours in the spectrum.
Moreover, rainbows have bands beyond red and violet in the respective near infrared and ultraviolet regions, however, these bands are not visible to humans. Only near frequencies of these regions to the visible spectrum are included in rainbows, since water and air become increasingly opaque to these frequencies, scattering the light. The UV band is sometimes visible to cameras using black and white film.
The question of whether everyone sees seven colours in a rainbow is related to the idea of linguistic relativity. Suggestions have been made that there is universality in the way that a rainbow is perceived. However, more recent research suggests that the number of distinct colours observed and what these are called depend on the language that one uses, with people whose language has fewer colour words seeing fewer discrete colour bands.
Explanation
When sunlight encounters a raindrop, part of the light is reflected and the rest enters the raindrop. The light is refracted at the surface of the raindrop. When this light hits the back of the raindrop, some of it is reflected off the back. When the internally reflected light reaches the surface again, once more some is internally reflected and some is refracted as it exits the drop. (The light that reflects off the drop, exits from the back, or continues to bounce around inside the drop after the second encounter with the surface, is not relevant to the formation of the primary rainbow.) The overall effect is that part of the incoming light is reflected back over the range of 0° to 42°, with the most intense light at 42°. This angle is independent of the size of the drop, but does depend on its refractive index. Seawater has a higher refractive index than rain water, so the radius of a "rainbow" in sea spray is smaller than that of a true rainbow. This is visible to the naked eye by a misalignment of these bows.
The reason the returning light is most intense at about 42° is that this is a turning point – light hitting the outermost ring of the drop gets returned at less than 42°, as does the light hitting the drop nearer to its centre. There is a circular band of light that all gets returned right around 42°. If the Sun were a laser emitting parallel, monochromatic rays, then the luminance (brightness) of the bow would tend toward infinity at this angle if interference effects are ignored . But since the Sun's luminance is finite and its rays are not all parallel (it covers about half a degree of the sky) the luminance does not go to infinity. Furthermore, the amount by which light is refracted depends upon its wavelength, and hence its colour. This effect is called dispersion. Blue light (shorter wavelength) is refracted at a greater angle than red light, but due to the reflection of light rays from the back of the droplet, the blue light emerges from the droplet at a smaller angle to the original incident white light ray than the red light. Due to this angle, blue is seen on the inside of the arc of the primary rainbow, and red on the outside. The result of this is not only to give different colours to different parts of the rainbow, but also to diminish the brightness. (A "rainbow" formed by droplets of a liquid with no dispersion would be white, but brighter than a normal rainbow.)
The light at the back of the raindrop does not undergo total internal reflection, and most of the light emerges from the back. However, light coming out the back of the raindrop does not create a rainbow between the observer and the Sun because spectra emitted from the back of the raindrop do not have a maximum of intensity, as the other visible rainbows do, and thus the colours blend together rather than forming a rainbow.
A rainbow does not exist at one particular location. Many rainbows exist; however, only one can be seen depending on the particular observer's viewpoint as droplets of light illuminated by the sun. All raindrops refract and reflect the sunlight in the same way, but only the light from some raindrops reaches the observer's eye. This light is what constitutes the rainbow for that observer. The whole system composed by the Sun's rays, the observer's head, and the (spherical) water drops has an axial symmetry around the axis through the observer's head and parallel to the Sun's rays. The rainbow is curved because the set of all the raindrops that have the right angle between the observer, the drop, and the Sun, lie on a cone pointing at the sun with the observer at the tip. The base of the cone forms a circle at an angle of 40–42° to the line between the observer's head and their shadow but 50% or more of the circle is below the horizon, unless the observer is sufficiently far above the earth's surface to see it all, for example in an aeroplane (see below). Alternatively, an observer with the right vantage point may see the full circle in a fountain or waterfall spray. Conversely, at lower latitudes near midday (specifically, when the sun's elevation exceeds 42 degrees) a rainbow will not be visible against the sky.
Mathematical derivation
It is possible to determine the perceived angle which the rainbow subtends as follows.
Given a spherical raindrop, and defining the perceived angle of the rainbow as , and the angle of the internal reflection as , then the angle of incidence of the Sun's rays with respect to the drop's surface normal is . Since the angle of refraction is , Snell's law gives us
,
where is the refractive index of water. Solving for , we get
.
The rainbow will occur where the angle is maximum with respect to the angle . Therefore, from calculus, we can set , and solve for , which yields
Substituting back into the earlier equation for yields ≈ 42° as the radius angle of the rainbow.
For red light (wavelength 750nm, based on the dispersion relation of water), the radius angle is 42.5°; for blue light (wavelength 350nm, ), the radius angle is 40.6°.
Variations
Double rainbows
A secondary rainbow, at a greater angle than the primary rainbow, is often visible. The term double rainbow is used when both the primary and secondary rainbows are visible. In theory, all rainbows are double rainbows, but since the secondary bow is always fainter than the primary, it may be too weak to spot in practice.
Secondary rainbows are caused by a double reflection of sunlight inside the water droplets. Technically the secondary bow is centred on the sun itself, but since its angular size is more than 90° (about 127° for violet to 130° for red), it is seen on the same side of the sky as the primary rainbow, about 10° outside it at an apparent angle of 50–53°. As a result of the "inside" of the secondary bow being "up" to the observer, the colours appear reversed compared to those of the primary bow.
The secondary rainbow is fainter than the primary because more light escapes from two reflections compared to one and because the rainbow itself is spread over a greater area of the sky. Each rainbow reflects white light inside its coloured bands, but that is "down" for the primary and "up" for the secondary. The dark area of unlit sky lying between the primary and secondary bows is called Alexander's band, after Alexander of Aphrodisias, who first described it.
Twinned rainbow
Unlike a double rainbow that consists of two separate and concentric rainbow arcs, the very rare twinned rainbow appears as two rainbow arcs that split from a single base. The colours in the second bow, rather than reversing as in a secondary rainbow, appear in the same order as the primary rainbow. A "normal" secondary rainbow may be present as well. Twinned rainbows can look similar to, but should not be confused with supernumerary bands. The two phenomena may be told apart by their difference in colour profile: supernumerary bands consist of subdued pastel hues (mainly pink, purple and green), while the twinned rainbow shows the same spectrum as a regular rainbow.
The cause of a twinned rainbow is believed to be the combination of different sizes of water drops falling from the sky. Due to air resistance, raindrops flatten as they fall, and flattening is more prominent in larger water drops. When two rain showers with different-sized raindrops combine, they each produce slightly different rainbows which may combine and form a twinned rainbow.
A numerical ray tracing study showed that a twinned rainbow on a photo could be explained by a mixture of 0.40 and 0.45 mm droplets. That small difference in droplet size resulted in a small difference in flattening of the droplet shape, and a large difference in flattening of the rainbow top.
Meanwhile, the even rarer case of a rainbow split into three branches was observed and photographed in nature.
Full-circle rainbow
In theory, every rainbow is a circle, but from the ground, usually only its upper half can be seen. Since the rainbow's centre is diametrically opposed to the Sun's position in the sky, more of the circle comes into view as the sun approaches the horizon, meaning that the largest section of the circle normally seen is about 50% during sunset or sunrise. Viewing the rainbow's lower half requires the presence of water droplets below the observer's horizon, as well as sunlight that is able to reach them. These requirements are not usually met when the viewer is at ground level, either because droplets are absent in the required position, or because the sunlight is obstructed by the landscape behind the observer. From a high viewpoint such as a high building or an aircraft, however, the requirements can be met and the full-circle rainbow can be seen. Like a partial rainbow, the circular rainbow can have a secondary bow or supernumerary bows as well. It is possible to produce the full circle when standing on the ground, for example by spraying a water mist from a garden hose while facing away from the sun.
A circular rainbow should not be confused with the glory, which is much smaller in diameter and is created by different optical processes. In the right circumstances, a glory and a (circular) rainbow or fog bow can occur together. Another atmospheric phenomenon that may be mistaken for a "circular rainbow" is the 22° halo, which is caused by ice crystals rather than liquid water droplets, and is located around the Sun (or Moon), not opposite it.
Supernumerary rainbows
In certain circumstances, one or several narrow, faintly coloured bands can be seen bordering the violet edge of a rainbow; i.e., inside the primary bow or, much more rarely, outside the secondary. These extra bands are called supernumerary rainbows or supernumerary bands; together with the rainbow itself the phenomenon is also known as a stacker rainbow. The supernumerary bows are slightly detached from the main bow, become successively fainter along with their distance from it, and have pastel colours (consisting mainly of pink, purple and green hues) rather than the usual spectrum pattern. The effect becomes apparent when water droplets are involved that have a diameter of about 1 mm or less; the smaller the droplets are, the broader the supernumerary bands become, and the less saturated their colours. Due to their origin in small droplets, supernumerary bands tend to be particularly prominent in fogbows.
Supernumerary rainbows cannot be explained using classical geometric optics. The alternating faint bands are caused by interference between rays of light following slightly different paths with slightly varying lengths within the raindrops. Some rays are in phase, reinforcing each other through constructive interference, creating a bright band; others are out of phase by up to half a wavelength, cancelling each other out through destructive interference, and creating a gap. Given the different angles of refraction for rays of different colours, the patterns of interference are slightly different for rays of different colours, so each bright band is differentiated in colour, creating a miniature rainbow. Supernumerary rainbows are clearest when raindrops are small and of uniform size. The very existence of supernumerary rainbows was historically a first indication of the wave nature of light, and the first explanation was provided by Thomas Young in 1804.
Reflected rainbow, reflection rainbow
When a rainbow appears above a body of water, two complementary mirror bows may be seen below and above the horizon, originating from different light paths. Their names are slightly different.
A reflected rainbow may appear in the water surface below the horizon. The sunlight is first deflected by the raindrops, and then reflected off the body of water, before reaching the observer. The reflected rainbow is frequently visible, at least partially, even in small puddles.
A reflection rainbow may be produced where sunlight reflects off a body of water before reaching the raindrops, if the water body is large, quiet over its entire surface, and close to the rain curtain. The reflection rainbow appears above the horizon. It intersects the normal rainbow at the horizon, and its arc reaches higher in the sky, with its centre as high above the horizon as the normal rainbow's centre is below it. Reflection bows are usually brightest when the sun is low because at that time its light is most strongly reflected from water surfaces. As the sun gets lower the normal and reflection bows are drawn closer together. Due to the combination of requirements, a reflection rainbow is rarely visible.
Up to eight separate bows may be distinguished if the reflected and reflection rainbows happen to occur simultaneously: the normal (non-reflection) primary and secondary bows above the horizon (1, 2) with their reflected counterparts below it (3, 4), and the reflection primary and secondary bows above the horizon (5, 6) with their reflected counterparts below it (7, 8).
Monochrome rainbow
Occasionally a shower may happen at sunrise or sunset, where the shorter wavelengths like blue and green have been scattered and essentially removed from the spectrum. Further scattering may occur due to the rain, and the result can be the rare and dramatic monochrome or red rainbow.
Higher-order rainbows
In addition to the common primary and secondary rainbows, it is also possible for rainbows of higher orders to form. The order of a rainbow is determined by the number of light reflections inside the water droplets that create it: One reflection results in the first-order or primary rainbow; two reflections create the second-order or secondary rainbow. More internal reflections cause bows of higher orders—theoretically unto infinity. As more and more light is lost with each internal reflection, however, each subsequent bow becomes progressively dimmer and therefore increasingly difficult to spot. An additional challenge in observing the third-order (or tertiary) and fourth-order (quaternary) rainbows is their location in the direction of the sun (about 40° and 45° from the sun, respectively), causing them to become drowned in its glare.
For these reasons, naturally occurring rainbows of an order higher than 2 are rarely visible to the naked eye. Nevertheless, sightings of the third-order bow in nature have been reported, and in 2011 it was photographed definitively for the first time. Shortly after, the fourth-order rainbow was photographed as well, and in 2014 the first ever pictures of the fifth-order (or quinary) rainbow were published. The quinary rainbow lies partially in the gap between the primary and secondary rainbows and is far fainter than even the secondary. In a laboratory setting, it is possible to create bows of much higher orders. Felix Billet (1808–1882) depicted angular positions up to the 19th-order rainbow, a pattern he called a "rose of rainbows". In the laboratory, it is possible to observe higher-order rainbows by using extremely bright and well collimated light produced by lasers. Up to the 200th-order rainbow was reported by Ng et al. in 1998 using a similar method but an argon ion laser beam.
Tertiary and quaternary rainbows should not be confused with "triple" and "quadruple" rainbows—terms sometimes erroneously used to refer to the (much more common) supernumerary bows and reflection rainbows.
Rainbows under moonlight
Like most atmospheric optical phenomena, rainbows can be caused by light from the Sun, but also from the Moon. In case of the latter, the rainbow is referred to as a lunar rainbow or moonbow. They are much dimmer and rarer than solar rainbows, requiring the Moon to be near-full in order for them to be seen. For the same reason, moonbows are often perceived as white and may be thought of as monochrome. The full spectrum is present, however, but the human eye is not normally sensitive enough to see the colours. Long exposure photographs will sometimes show the colour in this type of rainbow.
Fogbow
Fogbows form in the same way as rainbows, but they are formed by much smaller cloud and fog droplets that diffract light extensively. They are almost white with faint reds on the outside and blues inside; often one or more broad supernumerary bands can be discerned inside the inner edge. The colours are dim because the bow in each colour is very broad and the colours overlap. Fogbows are commonly seen over water when air in contact with the cooler water is chilled, but they can be found anywhere if the fog is thin enough for the sun to shine through and the sun is fairly bright. They are very large—almost as big as a rainbow and much broader. They sometimes appear with a glory at the bow's centre.
Fog bows should not be confused with ice halos, which are very common around the world and visible much more often than rainbows (of any order), yet are unrelated to rainbows.
Sleetbow
A sleetbow forms in the same way as a typical rainbow, with the exception that it occurs when light passes through falling sleet (ice pellets) instead of liquid water. As light passes through the sleet, the light is refracted causing the rare phenomena. These have been documented across United States with the earliest publicly documented and photographed sleetbow being seen in Richmond, Virginia on 21 December 2012. Just like regular rainbows, these can also come in various forms, with a monochrome sleetbow being documented on 7 January 2016 in Valparaiso, Indiana.
Circumhorizontal and circumzenithal arcs
The circumzenithal and circumhorizontal arcs are two related optical phenomena similar in appearance to a rainbow, but unlike the latter, their origin lies in light refraction through hexagonal ice crystals rather than liquid water droplets. This means that they are not rainbows, but members of the large family of halos.
Both arcs are brightly coloured ring segments centred on the zenith, but in different positions in the sky: The circumzenithal arc is notably curved and located high above the Sun (or Moon) with its convex side pointing downwards (creating the impression of an "upside down rainbow"); the circumhorizontal arc runs much closer to the horizon, is more straight and located at a significant distance below the Sun (or Moon). Both arcs have their red side pointing towards the Sun and their violet part away from it, meaning the circumzenithal arc is red on the bottom, while the circumhorizontal arc is red on top.
The circumhorizontal arc is sometimes referred to by the misnomer "fire rainbow". In order to view it, the Sun or Moon must be at least 58° above the horizon, making it a rare occurrence at higher latitudes. The circumzenithal arc, visible only at a solar or lunar elevation of less than 32°, is much more common, but often missed since it occurs almost directly overhead.
Extraterrestrial rainbows
It has been suggested that rainbows might exist on Saturn's moon Titan, as it has a wet surface and humid clouds. The radius of a Titan rainbow would be about 49° instead of 42°, because the fluid in that cold environment is methane instead of water. Although visible rainbows may be rare due to Titan's hazy skies, infrared rainbows may be more common, but an observer would need infrared night vision goggles to see them.
Rainbows with different materials
Droplets (or spheres) composed of materials with different refractive indices than plain water produce rainbows with different radius angles. Since salt water has a higher refractive index, a sea spray bow does not perfectly align with the ordinary rainbow, if seen at the same spot. Tiny plastic or glass marbles may be used in road marking as a reflectors to enhance its visibility by drivers at night. Due to a much higher refractive index, rainbows observed on such marbles have a noticeably smaller radius. One can easily reproduce such phenomena by sprinkling liquids of different refractive indices in the air, as illustrated in the photo.
The displacement of the rainbow due to different refractive indices can be pushed to a peculiar limit. For a material with a refractive index larger than 2, there is no angle fulfilling the requirements for the first order rainbow. For example, the index of refraction of diamond is about 2.4, so diamond spheres would produce rainbows starting from the second order, omitting the first order. In general, as the refractive index exceeds a number , where is a natural number, the critical incidence angle for times internally reflected rays escapes the domain . This results in a rainbow of the -th order shrinking to the antisolar point and vanishing.
Scientific history
The classical Greek scholar Aristotle (384–322 BC) was first to devote serious attention to the rainbow. According to Raymond L. Lee and Alistair B. Fraser, "Despite its many flaws and its appeal to Pythagorean numerology, Aristotle's qualitative explanation showed an inventiveness and relative consistency that was unmatched for centuries. After Aristotle's death, much rainbow theory consisted of reaction to his work, although not all of this was uncritical."
In Book I of Naturales Quaestiones (), the Roman philosopher Seneca the Younger discusses various theories of the formation of rainbows extensively, including those of Aristotle. He notices that rainbows appear always opposite to the Sun, that they appear in water sprayed by a rower, in the water spat by a fuller on clothes stretched on pegs or by water sprayed through a small hole in a burst pipe. He even speaks of rainbows produced by small rods (virgulae) of glass, anticipating Newton's experiences with prisms. He takes into account two theories: one, that the rainbow is produced by the Sun reflecting in each water drop, the other, that it is produced by the Sun reflected in a cloud shaped like a concave mirror; he favours the latter. He also discusses other phenomena related to rainbows: the mysterious "virgae" (rods), halos and parhelia.
According to Hüseyin Gazi Topdemir, the Arab physicist and polymath Ibn al-Haytham (965–1039 AD) attempted to provide a scientific explanation for the rainbow phenomenon. In his Maqala fi al-Hala wa Qaws Quzah (On the Rainbow and Halo), al-Haytham "explained the formation of rainbow as an image, which forms at a concave mirror. If the rays of light coming from a farther light source reflect to any point on axis of the concave mirror, they form concentric circles in that point. When it is supposed that the sun as a farther light source, the eye of viewer as a point on the axis of mirror and a cloud as a reflecting surface, then it can be observed the concentric circles are forming on the axis." He was not able to verify this because his theory that "light from the sun is reflected by a cloud before reaching the eye" did not allow for a possible experimental verification. This explanation was repeated by Averroes, and, though incorrect, provided the groundwork for the correct explanations later given by Kamāl al-Dīn al-Fārisī in 1309 and, independently, by Theodoric of Freiberg (c. 1250–c. 1311)—both having studied al-Haytham's Book of Optics.
In Song dynasty China (960–1279), a polymath scholar-official named Shen Kuo (1031–1095) hypothesised—as a certain Sun Sikong (1015–1076) did before him—that rainbows were formed by a phenomenon of sunlight encountering droplets of rain in the air. Paul Dong writes that Shen's explanation of the rainbow as a phenomenon of atmospheric refraction "is basically in accord with modern scientific principles."
According to Nader El-Bizri, the Persian astronomer, Qutb al-Din al-Shirazi (1236–1311), gave a fairly accurate explanation for the rainbow phenomenon. This was elaborated on by his student, Kamāl al-Dīn al-Fārisī (1267–1319), who gave a more mathematically satisfactory explanation of the rainbow. He "proposed a model where the ray of light from the sun was refracted twice by a water droplet, one or more reflections occurring between the two refractions." An experiment with a water-filled glass sphere was conducted and al-Farisi showed the additional refractions due to the glass could be ignored in his model. As he noted in his Kitab Tanqih al-Manazir (The Revision of the Optics), al-Farisi used a large clear vessel of glass in the shape of a sphere, which was filled with water, in order to have an experimental large-scale model of a rain drop. He then placed this model within a camera obscura that has a controlled aperture for the introduction of light. He projected light unto the sphere and ultimately deduced through several trials and detailed observations of reflections and refractions of light that the colours of the rainbow are phenomena of the decomposition of light.
In Europe, Ibn al-Haytham's Book of Optics was translated into Latin and studied by Robert Grosseteste. His work on light was continued by Roger Bacon, who wrote in his of 1268 about experiments with light shining through crystals and water droplets showing the colours of the rainbow. In addition, Bacon was the first to calculate the angular size of the rainbow. He stated that the rainbow summit can not appear higher than 42° above the horizon. Theodoric of Freiberg is known to have given an accurate theoretical explanation of both the primary and secondary rainbows in 1307. He explained the primary rainbow, noting that "when sunlight falls on individual drops of moisture, the rays undergo two refractions (upon ingress and egress) and one reflection (at the back of the drop) before transmission into the eye of the observer." He explained the secondary rainbow through a similar analysis involving two refractions and two reflections.
Descartes' 1637 treatise, Discourse on Method, further advanced this explanation. Knowing that the size of raindrops did not appear to affect the observed rainbow, he experimented with passing rays of light through a large glass sphere filled with water. By measuring the angles that the rays emerged, he concluded that the primary bow was caused by a single internal reflection inside the raindrop and that a secondary bow could be caused by two internal reflections. He supported this conclusion with a derivation of the law of refraction (subsequently to, but independently of, Snell) and correctly calculated the angles for both bows. His explanation of the colours, however, was based on a mechanical version of the traditional theory that colours were produced by a modification of white light.
Isaac Newton demonstrated that white light was composed of the light of all the colours of the rainbow, which a glass prism could separate into the full spectrum of colours, rejecting the theory that the colours were produced by a modification of white light. He also showed that red light is refracted less than blue light, which led to the first scientific explanation of the major features of the rainbow. Newton's corpuscular theory of light was unable to explain supernumerary rainbows, and a satisfactory explanation was not found until Thomas Young realised that light behaves as a wave under certain conditions, and can interfere with itself.
Young's work was refined in the 1820s by George Biddell Airy, who explained the dependence of the strength of the colours of the rainbow on the size of the water droplets. Modern physical descriptions of the rainbow are based on Mie scattering, work published by Gustav Mie in 1908. Advances in computational methods and optical theory continue to lead to a fuller understanding of rainbows. For example, Nussenzveig provides a modern overview.
Experiments
Experiments on the rainbow phenomenon using artificial raindrops, i.e. water-filled spherical flasks, go back at least to Theodoric of Freiberg in the 14th century. Later, also Descartes studied the phenomenon using a Florence flask. A flask experiment known as Florence's rainbow is still often used today as an imposing and intuitively accessible demonstration experiment of the rainbow phenomenon. It consists in illuminating (with parallel white light) a water-filled spherical flask through a hole in a screen. A rainbow will then appear thrown back / projected on the screen, provided the screen is large enough. Due to the finite wall thickness and the macroscopic character of the artificial raindrop, several subtle differences exist as compared to the natural phenomenon, including slightly changed rainbow angles and a splitting of the rainbow orders.
A very similar experiment consists in using a cylindrical glass vessel filled with water or a solid transparent cylinder and illuminated either parallel to the circular base (i.e. light rays remaining at a fixed height while they transit the cylinder) or under an angle to the base. Under these latter conditions the rainbow angles change relative to the natural phenomenon since the effective index of refraction of water changes (Bravais' index of refraction for inclined rays applies).
Other experiments use small liquid drops (see text above).
Culture and mythology
Rainbows occur frequently in mythology, and have been used in the arts. The first literary occurrence of a rainbow is in the Book of Genesis chapter 9, as part of the flood story of Noah, where it is a sign of God's covenant to never destroy all life on Earth with a global flood again. In Norse mythology, the rainbow bridge Bifröst connects the world of men (Midgard) and the realm of the gods (Asgard). Cuchavira was the god of the rainbow for the Muisca in present-day Colombia and when the regular rains on the Bogotá savanna were over the people thanked him, offering gold, snails and small emeralds. Some forms of Tibetan Buddhism or Dzogchen reference a rainbow body. The Irish leprechaun's secret hiding place for his pot of gold is usually said to be at the end of the rainbow. This place is appropriately impossible to reach, because the rainbow is an optical effect which cannot be approached. In Greek mythology, the goddess Iris is the personification of the rainbow, a messenger goddess who, like the rainbow, connects the mortal world with the gods through messages. In Albanian folk beliefs the rainbow is regarded as the belt of the goddess Prende, and oral legend has it that anyone who jumps over the rainbow changes their sex.
In heraldry, the rainbow proper consists of 4 bands of colour (argent, gules, or, and vert) with the ends resting on clouds. Generalised examples in coat of arms include those of the towns of Regen and Pfreimd, both in Bavaria, Germany; of Bouffémont, France; and of the 69th Infantry Regiment (New York) of the United States Army National Guard.
Rainbow flags have been used for centuries. It was a symbol of the Cooperative movement in the German Peasants' War in the 16th century, of peace in Italy, and of LGBT pride and LGBT social movements; the rainbow flag as a symbol of LGBT pride and the June pride month since it was designed by Gilbert Baker in 1978. In 1994, Archbishop Desmond Tutu and President Nelson Mandela described newly democratic post-apartheid South Africa as the rainbow nation. The rainbow has also been used in technology product logos, including the Apple computer logo. Many political alliances spanning multiple political parties have called themselves a "Rainbow Coalition".
Pointing at rainbows has been considered a taboo in many cultures.
In Saudi Arabia and other similar-minded countries, authorities seize children's clothing (including hats, hair clips, pencil cases, etc.) and toys if they are rainbow-coloured, claiming that such can encourage homosexuality, and selling such items is illegal.
| Physical sciences | Atmospheric optics | null |
2072335 | https://en.wikipedia.org/wiki/Lambeosaurus | Lambeosaurus | Lambeosaurus ( , meaning "Lambe's lizard") is a genus of hadrosaurid dinosaur that lived about 75 million years ago, in the Late Cretaceous period (Campanian stage) of North America. This bipedal/quadrupedal, herbivorous dinosaur is known for its distinctive hollow cranial crest, which in the best-known species resembled a mitten. Several possible species have been named, from Canada, the United States, and Mexico, but only the two Canadian species are currently recognized as valid.
Material relevant to the genus was first named by Lawrence Lambe in 1902. Over twenty years later, the modern name was coined in 1923 by William Parks, in honour of Lambe, based on better preserved specimens. The genus has a complicated taxonomic history, in part because small-bodied crested hadrosaurids now recognized as juveniles were once thought to belong to their own genera and species. Currently, the various skulls assigned to the type species L. lambei are interpreted as showing age differences and sexual dimorphism. Lambeosaurus was closely related to the better known Corythosaurus, which is found in slightly older rocks, as well as the less well-known genera Hypacrosaurus and Olorotitan. All had unusual crests, which are now generally assumed to have served social functions like noisemaking and recognition.
Discovery and species
Naming of Lambeosaurus
In the 1880s and 1890s, expeditions of the Geological Survey of Canada into Alberta identified that the rocks along the Red Deer River bore dinosaurs of scientific importance. These deposits were identified as either the Edmonton or Belly River Series, of the middle to end Cretaceous. Canadian palaeontologist Lawrence M. Lambe participated in three expeditions in 1897, 1898 and 1901, narrowing down a locality along an extensive series of badlands between Berry Creek and Dear Lodge canyon as the locality for excavations. The greatest difficulty with the dinosaur bones discovered was that they had been scattered before burial and very fragile, but from these expeditions Lambe acquired material of Deinodon, Ornithomimus, Palaeoscincus, Monoclonius, the new armoured dinosaurs Stereocephalus and Stegoceras, and three new species of Trachodon of the subgenus Pteropelyx that he named Trachodon selwyni, Trachodon marginatus, and Trachodon altidens in 1902. T. selwyni, named for former Geological Survey director Alfred Selwyn, was established for a jaw with many teeth, and a provisionally-referred femur, T. marginatus was named for the partial skeleton of one individual, as well as isolated jaw and limb bones, and T. altidens was named for a partial maxilla with many teeth. American palaeontologist Henry Fairfield Osborn, summarizing the fauna of the mid-Cretaceous across all of North America in the preceding article of the same publication, briefly discussed Lambe's species, and even provided the possible new subgenus name Didanodon for T. altidens.
A later Geological Survey expedition in 1913 resulted in American paleontologist Charles Hazelius Sternberg discovering a specimen of hadrosaur with skin impressions from along the same stretch of river, which Lambe described in 1914 and referred to Trachodon marginatus. Sternberg also found a second specimen in the same expedition, and the two specimens combined showed that the skin and pelvis of T. marginatus could distinguish it from other hadrosaurs, including a specimen that had been provisionally referred from the younger Edmonton series. Lambe described the specimens more thoroughly later in 1914, and determined from the nearly complete skull of one individual that T. marginatus belonged in a new genus, for which he named Stephanosaurus. To Stephanosaurus he referred the original jaws, skeleton, and isolated elements he described in 1902, as well as the two specimens found in 1913 by Sternberg. However, American palaeontologist Barnum Brown highlighted later that year that the referred specimens of Stephanosaurus may not belong to the same taxon as the type skeleton, which preserves the partial articulated forelimb, but not any material of the skull. As a result, Brown considered the identity of the skulls as uncertain, highlighting similarities with his new genus Corythosaurus. Lambe, in 1920, referred another even more perfect skull to Stephanosaurus, showing separation from Corythosaurus. However, Canadian palaeontologist William A. Parks noted that the logic of Brown still applied, and designated the new genus and species Lambeosaurus lambei for the complete skulls that could not be justifiably referred to Stephanosaurus, to give Lambe (who was deceased) as much credit as possible for the new hadrosaur.
American palaeontologist Charles W. Gilmore more completely described the type material of Lambeosaurus and provided discussion of Stephanosaurus in 1924. He found that the type material of Stephanosaurus, Canadian Museum of Nature specimen 419, could not be the same taxon as the skulls. The better-preserved skull that Lambe described in 1920, CMN 2869, was identified as the type specimen, having been found by Charles Mortram Sternberg in 1917 from around southeast of mouth of Little Sandhill Creek, below the top of the Belly River beds. Additional referred specimens include CMN 351, a half skull with limb bones that is one of the specimens found by C.H. Sternberg in 1913 southeast of the mouth of Berry Creek, and CMN 8503, a partial skull and articulated skeleton also found by C.M. Sternberg in 1917 west of the mouth of Little Sandhill Creek.
Procheneosaurus and Tetragonosaurus
The American Museum of Natural History also excavated in the region of the Red Deer River, primarily through the work of Brown from 1909 to 1914. One individual discovered, prepared into a panel mount for display at the AMNH, was referred to by the name Procheneosaurus by American paleontologist William Diller Matthew in 1920, as a particularly small hadrosaur. In 1931, Parks described more small crested hadrosaurs found by expeditions of the University of Toronto into the Red Deer badlands. The first was named for the Royal Ontario Museum specimen 3577, which was found on above the river level and southeast of Sand Creek, being excavated and mounted by American paleontologist Levi Sternberg. Parks named this specimen, which included a skull and some of the vertebral column, as the new taxon Tetragonosaurus praeceps. The second specimen, ROM 3578, was found in 1927 by Levi Sternberg from above the river level and downriver of Sand Creek. This specimen, including only the skull, was named Tetragonosaurus erectofrons by Parks. Together, both species were considered close to Cheneosaurus, and showed similar low, domed crests.
Charles M. Sternberg identified that the crested hadrosaurs were in need of reassessment, and began his 1935 study of them by finishing the preparation of the numerous skulls and skeletons present at the Canadian Museum of Nature. Among these 18 specimens, Sternberg found that CMN 8703, a nearly complete skeleton with skin, could be referred to Lambeosaurus, while CMN 8503, previously referred, was instead a specimen of Corythosaurus. A second undescribed species of Lambeosaurus was identified from a skull and skeleton, as well as a new species of Tetragonosaurus. This new species of Tetragonosaurus was named T. cranibrevis, based on the partial skull CMN 8633 found by C.M. Sternberg in 1928 at a level of above the river and south of the mouth of Berry Creek. The new species of Lambeosaurus was named L. clavinitialis for the skull and skeleton CMN 8703 found by C.M. Sternberg in 1928 at a level of above the river and south of the mouth of Berry Creek. Finally, Sternberg named the new species Lambeosaurus magnicristatum for the mostly complete skull and skeleton CMN 8705 found by C.M. Sternberg in 1919 southwest of the mouth of Little Sandhill Creek and near the top of the beds. Parks also named the new species Corythosaurus frontalis for ROM 869 from the same area.
The extensive species and genera of hadrosaurs was reviewed in 1942 by American paleontologists Richard Swann Lull and Nelda E. Wright. Lull and Wright considered L. lambei, L. clavinitialis, and L. magnicristatus as valid and distinguishable species of Lambeosaurus, each known from multiple specimens. To L. lambei they referred CMN 2869, ROM 5131, and ROM 1218, the second collected in 1920 by Levi Sternberg northeast of Happy Jack Ferry, while the latter was collected in the 1919 expedition on the south side of Red Deer River. L. clavinitialis included the specimens CMN 8703, CMN 351, YPM 3222, the latter found in 1919 by C.M. Sternberg south of the mouth of Little Sandhill Creek. L. magnicristatus was limited to the type CMN 8705, but numerous specimens could be referred to Lambeosaurus but not a species within (CMN 8502, AMNH 5353, AMNH 5373, AMNH 5666, and USNM 10309).
Lull and Wright interpreted the description of Procheneosaurus as sufficient to identify that Matthew was referring to AMNH 5340, a nearly complete skull and skeleton, that could not be distinguished from the Tetragonosaurus species on a generic level. They interpreted that Parks naming Tetragonosaurus was under the belief that Procheneosaurus was invalidly described, and as Procheneosaurus and Tetragonosaurus were synonymous with the former lacking a species, they established Procheneosaurus praeceps as the type species. As a result, Lull and Wright also created the combinations Procheneosaurus erectofrons and Procheneosaurus cranibrevis, for the other former species of Tetragonosaurus. Trachodon altidens was also referred to Procheneosaurus, as P. altidens. This treatment was upheld by a petition by Lull to the International Commission on Zoological Nomenclature, that ruled in 1947 in favour of Procheneosaurus as a valid name with seniority over Tetragonosaurus and that P. praeceps would be the type species by including AMNH 5340. The commission also considered that Tetragonosaurus was not an available name due to lacking a clearly designated type species. In addition to their type specimens (ROM 3577 for P. praeceps, ROM 3578 for P. erectofrons, CMN 8633 for P. cranibrevis and CMN 1092 for P. altidens), Lull and Wright referred AMNH 5340 to P. praeceps, AMNH 5461 and AMNH 5469 to P. erectofrons, with the latter two specimens collected by Brown from Montana in 1916 and including a skull and much of a skeleton, and a partial skeleton, respectively.
The hadrosaur taxonomy of Lull and Wright was followed by American palaeontologist John Ostrom who described the skulls of the different Lambeosaurus and Procheneosaurus species, as well as reassessed the fragmentary species Hadrosaurus paucidens, considering it particularly similar to Lambeosaurus lambei and reassigning it as Lambeosaurus paucidens. The use of Procheneosaurus was furthered by Russian palaeontologist Anatoly K. Rozhdestvensky, who described in 1968 some new hadrosaurs that linked between the Asian and North American faunas. One of these new hadrosaurs was the species he named Procheneosaurus convincens, for an almost complete skeleton and skull accessioned at the Paleontologital Institute as PIN 2230. This skeleton had been found in 1961 at the Shakh-Shakh locality north of Tashkent, as the most complete dinosaur discovered in Kazakhstan, and came from the Santonian aged Dabrazin Formation.
Identification of juveniles
In 1975, American palaeontologist Peter Dodson published a study assessing the impacts of growth on the anatomy of lambeosaurine hadrosaurs. Based on the changes shown by living reptiles during growth, Dodson reassessed the three genera and 12 species of lambeosaur believed to have lived in the Oldman Formation (historic Belly River series), which showed extreme cranial variation in their crests alone. He measured and plotted morphometrics of the crests of 36 individuals, including the types of all three Lambeosaurus species and all three North American Procheneosaurus species, and compared them with the sizes of the individuals, assuming that smaller individuals would be younger. From this, it could be seem that the smaller Procheneosaurus specimens appeared to form a gradual increase in size and crest form into the larger Lambeosaurus and Corythosaurus specimens, and were best considered juveniles. Dodson concluded that Procheneosaurus praeceps were juveniles of Lambeosaurus lambei, and Procheneosaurus erectofrons and P. cranibrevis were juveniles of Corythosaurus. Lambeosaurus magnicristatus remained distinct. Dodson also proposed that variation among individuals of the same size was due to sexual dimorphism, resulting in the identifications of CMN 8503 (referred to C. intermedius), ROM 869 (type of C. frontalis), CMN 351 and YPM 3222 (referred to L. clavinitialis), and CMN 8703 (type of L. clavinitialis) as females of L. lambei, and CMN 2869, ROM 1218, ROM 5131, AMNH 5353 and AMNH 5373 (all previously referred to L. lambei) as males of L. lambei. L. magnicristatus was represented by one male and one female individual, the type CMN 8705 and Royal Tyrrell Museum specimen 1966.04.1 respectively.
The interpretation of Dodson of the low-crested cheneosaurs as juveniles of Lambeosaurus, Corythosaurus, and Hypacrosaurus was followed by subsequent studies, with some slight adjustments. American paleontologist James A. Hopson suggested that in addition to Dodson's identification of the species of Procheneosaurus, L. clavinitialis could represent female individuals of L. magnicristatus rather than L. lambei. Polish palaeontologists Teresa Maryańska and Halszka Osmólska also identified that P. convincens is likely a separate species from the others and identified it as "Procheneosaurus" convincens. American paleontologist William J. Morris even noted that there is very little that separates Lambeosaurus, Corythosaurus, and Hypacrosaurus beyond the skull anatomy, and that a specimen from Baja California cannot be identified as either Hypacrosaurus or Lambeosaurus. These Baja California specimens were excavated from 1968 to 1974 in the El Gallo formation, and were more completely described from additional material by Morris in 1981, where he found that they were a new species he named ?Lambeosaurus laticaudus with the type specimen of Natural History Museum of Los Angeles County number 17715, only tentatively referring the species to Lambeosaurus. As well as material from Baja California assigned to the new species, material was described from the Bearpaw Formation of Montana that was tentatively referred to L. magnicristatus. This would represent the first lambeosaur specimen from marine sediments, though it only included partial jaw bones.
In the 1990 review of the Hadrosauridae by American palaeontologists David B. Weishampel and John R. Horner, Lambeosaurus was considered to include L. lambei, L. magnicristatus, and ?L. laticaudus, with Tetragonosaurus praeceps, L. clavinitialis, and Corythosaurus frontalis as synonyms of L. lambei. They also considered Didanodon and Procheneosaurus as synonymic genera with Lambeosaurus, while the other species described by Lambe: T. selwyni, T. altidens and T. marginatus, and the species Hadrosaurus paucidens, were dubious hadrosaurids that could not be classified. Procheneosaurus convincens was a juvenile and synonym of Jaxartosaurus aralensis from a similar locality, and T. erectofrons and T. cranibrevis were juveniles of Corythosaurus. American paleontologists David B. Norman and Hans-Dieter Sues conversely argued that Jaxartosaurus was too separate in time and space from Procheneosaurus convincens, and that though its validity was questionable, it could not be referred to any known genera.
Redescriptions of species
The taxonomy of Weishampel and Horner was reiterated in the 2004 review of Hadrosauridae by Horner and colleagues, though the taxonomy of the species of Tetragonosaurus would be revisited with a more detailed description of T. erectofrons in 2005 by Canadian paleontologist David C. Evans and colleagues. Dodson had not included the type specimen of T. erectofrons in his analysis of lambeosaurs, due to its partially incomplete skull. However, Evans and colleagues were able to identify features that separated Corythosaurus from Lambeosaurus regardless of age, therefore providing more completely justified assessments of the juveniles. They were able to identify diagnostic traits of Corythosaurus in T. erectofrons and a referred specimen of T. cranibrevis, but the type of T. cranibrevis showed the anatomy of Lambeosaurus making the species a synonym of it rather than Corythosaurus.
Following the descriptions of skulls and skeletons of various ontogenetic stages of L. lambei, L. magnicristatus remained the only lambeosaurine from the Cretaceous of Alberta to lack a description of its skeleton. Evans and Canadian paleontologist Robert R. Reisz redescribed L. magnicristatus in 2007, based primarily on TMP 1966.04.1. The type, CMN 8705, was originally a largely complete skeleton and skull when excavated, but prior to being named it was significantly damaged by water, destroying much of the limbs and girdles which were discarded, before it was named by Sternberg. TMP 1966.04.1 was originally discovered by C.M. Sternberg in 1937, around southeast of Manyberries, Alberta, and around below the top of the formation. It was originally excavated for the Canadian Museum of Nature, but was given to the Provincial Museum and Archives of Alberta in 1966, which later became the Royal Tyrrell Museum. When first discovered, the left side was exposed and had been badly weathered, probably over hundreds of years. To preserve it, Sternberg reinforced the skeleton with plaster and then jacketed it into five separate blocks that were shipped in straw-packed wooden crates to Ottawa. After being acquired by the PMAA, it was prepared for exhibition, with the better-preserved right side being exposed and bolted into a large wooden frame with styrofoam blocks cut out to hold it in place. It has been displayed since mid-1969. The redescription by Evans and Reisz showed that L. magnicristatus could be separated from L. lambei by anatomy as well as being a younger species. It also raised the question of whether L. clavinitialis were only individuals of L. lambei, or possibly belonged to a similar but separate species, though this required further study.
The suggestion of L. clavinitialis as a separate species from L. lambei was also considered by Evans in other studies, with new stratigraphic data suggesting that there is an age separation between the different "morphs", with temporal separation between L. lambei and L. clavinitialis, though they could still represent chronospecies or evolutionary lineages. As a result, L. clavinitialis has been identified as separate in following studies, occurring alongside L. lambei in the older deposits of the Dinosaur Park Formation, while L. magnicristatus is younger, but more restricted than L. lambei only being found within the lower formation. AMNH 5382, ROM 869, CMN 8703, YPM 3222 and TMP 1981.37.1 have all been referred to L. clavinitialis, while AMNH 5353, AMNH 5373, CMN 351, CMN 2869, CMN 8503, ROM 794, ROM 1218, FMNH 380, FMNH 1479, TMP 1982.38.1, and TMP 1997.012.0128 have been referred to L. lambei, and some specimens cannot be identified to the species level. It was included in the phylogenetic analysis of Hai Xing and colleagues in 2022 as a separate species for the first time and found to be more distant, with L. lambei and L. magnicristatus as sister taxa, and its crest morphometrics also differ significantly from L. lambei and L. magnicristatus.
In 2012, Spanish palaeontologist Albert Prieto-Márquez and colleagues redescribed the material of ?L. laticaudus, which was found to be closest to Velafrons, the only other lambeosaurine known from Mexico at the time. It was more distant from the species of Lambeosaurus, and showed enough anatomical differences that they gave it the new genus name Magnapaulia after Paul G. Haaga Jr. for his support of the Natural History Museum of Los Angeles County. Similarly, the skull of Procheneosaurus convincens was redescribed by Canadian palaeontologists Phil R. Bell and Kirstin S. Brink, who found that it could be differentiated from other lambeosaurines of similar age, and the taxa Jaxartosaurus and Aralosaurus found nearby, giving if the new genus name Kazaklambia in honour of the country it was found in. It has been found to be an early member of Lambeosaurini and relatively far from Lambeosaurus.
Description
Lambeosaurus, Corythosaurus and Hypacrosaurus are considered very close relatives of each other, and differ very little in anatomy beyond the skull and crest. They were large hadrosaurids, with highly developed jaws full of grinding teeth, a long tail stiffened by ossified tendons that prevented it from drooping, and more elongate limbs suggesting they were semi-quadrupedal and could move on both two legs and all fours as shown by footprints of related animals. The hands had four fingers, lacking the thumb of the generalized five-fingered tetrapod hand, and while the second, third, and fourth fingers were bunched together, the little finger was free and could be used to manipulate objects. Each foot had only the three central toes. L. lambei, L. magnicristatus , and L. clavinitialis would have reached around in length, and in weight, respectively. This is comparable to all species of Corythosaurus and Hypacrosaurus which were around in length and , and makes them the largest lambeosaurines except for Magnapaulia. Lambeosaurus is also one of many hadrosaurs to preserve the impressions of skin, which has been found across the neck, pelvis, legs, and tail.
Skull
Lambeosaurus was quite similar to Corythosaurus in everything but the form of the head adornment. Compared to Corythosaurus, the crest of Lambeosaurus, largely formed by the premaxillae, was shifted forward, and the hollow nasal passages within were at the front of the crest and stacked vertically. It also can be differentiated from Corythosaurus by its lack of forking nasal processes making up part of the sides of the crest, which is the only way to tell juveniles of the two genera apart, as the crests took on their distinctive forms as the animals aged.
The most distinctive feature, the crest, was different in the two well-known species. In L. lambei, it had a hatchet-like shape when the dinosaur was full-grown, and was somewhat shorter and more rounded in specimens interpreted as females. The "hatchet blade" projected in front of the eyes, and the "handle" was a solid bony rod that jutted out over the back of the skull. The "hatchet blade" had two sections: the uppermost portion was a thin bony "coxcomb" that grew out relatively late in life, when an individual neared adulthood; and the lower portion held hollow spaces that were continuations of the nasal passages. In L. magnicristatus, the "handle" was greatly reduced, and the "blade" expanded, forming a tall, exaggerated pompadour-like crest. This crest is damaged in the best overall specimen, and only the front half remains.
Skeleton
Lambeosaurus is known from several completely but briefly described skeletons, with no features of the skeleton distinguishing it from relatives, though L. lambei and L. magnicristatus can be separated by the skull and one feature of the pelvis. The neck of Lambeosaurus bears 14 or 15 cervical vertebrae, which fewer than Olorotitan but more than Parasaurolophus. They are generally consistent in form along the neck, with the exception of the first two cervicals which are specialized for supporting the skull with a taller above the vertebral body. The spines of the following cervicals are nearly absent, though the articulations between vertebrae are strong and (concave-convex) as in other hadrosaurids. The cervical vertebrae are very consistent in length, only varying between in L. magnicristatus specimen TMP 1966.04.1. There were 15 to 16 dorsal vertebrae present in the torso, with the spines much taller than in the cervicals and rectangular. The vertebrae of the tail have hexagonal articular faces, as in other hadrosaurs.
The in the shoulder girdle is an elongate, flat bone, long in TMP 1966.04.1 and gently curved as in Corythosaurus. The surface of the bone is relatively smooth, except for a large crest near the shoulder joint that serves to anchor muscles in the region. The is in the shape of a hatchet as in other hadrosaurids. The is relatively shorter than in Corythosaurus, but is still the most massive bone of the arm. It bears a large crest for the deltopectoral muscles that extends for half the length of the bone, before sharply merging into the shaft of the bone. The and are longer than the humerus as in other lambeosaurines, but are more robust than in Corythosaurus. While the ulna is one third longer than the humerus, it has very little expansion at the elbow or wrist. The hand is also more robust than in Corythosaurus, with longer digits relative to the , making it similar to Parasaurolophus. The second digit is the longest of the hand, despite the third and fourth metacarpals being longer, and the second and third digits bore hooves, which would have faced slightly inwards when walking.
Pelvic material of Lambeosaurus has been suggested to be different from other hadrosaurs, but the variation within Hypacrosaurus and Maiasaura show that these differences are probably individual and not related to species. The is elongate, with a humped upper margin that bears a prominent shelf overhanging the hip joint. Multiple scars from muscle attachments can be seen on the surface, and the (rear) process of the ilium is very similar between L. lambei, L. magnicristatus and Corythosaurus. The however, differs between the taxa: while in L. lambei and Corythosaurus it is bulbous with a large expansion in front of the hip joint, in L. magnicristatus this expansion is much smaller, and the entire projection is relatively shorter. The is subequal in length to the , with a slightly sigmoid outline and more massive proportions than Corythosaurus. The shaft of the bone is straight, but there are expansions at either end. Near the hip joint, the ischium broadens to articulate with the pubis and ilium, and at the underside of this region of the pelvis the ischium is notched. The opposite end of the ischium is sharply expanded into a pendant foot, which, though large and unique to lambeosaurines, is similar between Lambeosaurus species as being smaller than Parasaurolophus and Hypacrosaurus.
The femur is massive and columnar, and as in other hadrosaurs it is slightly longer than the of the lower leg. It is broad, with deep ridges for muscle attachments including a strongly developed and semicircular . The for the knee articulation are expanded enough to fully enclose a tunnel for extensor ligaments, giving a long articular surface. The tibia is massive and does not differ from other hadrosaurids, with the upper third of its length taken up by the that forms an arc to brace the from the front. The fibula is slender and the same length as the tibia, though its robustness is more similar to Hypacrosaurus than Corythosaurus. The femur is long in L. clavinitialis, while the humerus is and the ilium is . The foot in hadrosaurids is reduced to only three digits, which each bear spade-shaped hooves.
Integument
Impressions of the scales are known for several specimens; a specimen now assigned to L. lambei had a thin skin with uniform, polygonal scutes distributed in no particular order on the neck, torso, and tail. Similar scalation is known from the neck, forelimb, and foot of a specimen of L. magnicristatus.
Classification
Lambeosaurus is the type genus of the Lambeosaurinae, the subfamily of hadrosaurids that had hollow skull crests. Among the lambeosaurines, it is closely related to similar dinosaurs such as Corythosaurus and Hypacrosaurus, with little separating them but crest form. The relationships among these dinosaur genera are difficult to pick out. Some early classifications placed these genera in the tribe Corythosaurini, which was found by David Evans and Robert Reisz to include Lambeosaurus as the sister taxon to a clade made up of Corythosaurus, Hypacrosaurus, and the Russian genus Olorotitan; these lambeosaurines, with Nipponosaurus. However, later researchers pointed out that due to the rules of priority set forth by the ICZN, any tribe containing Lambeosaurus is properly named Lambeosaurini, and that therefore the name "Corythosaurini" is a junior synonym. The following cladogram illustrating the relationships of Lambeosaurus and its close relatives was recovered in a 2022 phylogenetic analysis by Xing Hai and colleagues, finding it to be a close relative of Amurosaurus.
Paleobiology
Feeding
As a hadrosaurid, Lambeosaurus was a large bipedal/quadrupedal herbivore, eating plants with a sophisticated skull that permitted a grinding motion analogous to mammalian chewing. Its teeth were continually replaced and were packed into dental batteries that each contained hundreds of teeth, only a relative handful of which were in use at any time. It used its beak to crop plant material, which was held in the jaws by a cheek-like organ. Feeding would have been from the ground up to around above. As noted by Bob Bakker, lambeosaurines have narrower beaks than hadrosaurines, implying that Lambeosaurus and its relatives could feed more selectively than their broad-beaked, crestless counterparts.
Cranial crest
Like other lambeosaurines such as Parasaurolophus and Corythosaurus, Lambeosaurus had a distinctive crest on the top of its head. Respiratory tracts, ending in a nasal cavity, ran back through this crest, making it mostly hollow. Many suggestions have been made for the function or functions of the crest, including housing salt glands, improving the sense of smell, use as a snorkel or air trap, acting as a resonating chamber for making sounds, or being a method for different species or different sexes of the same species to recognize each other. Social functions such as noisemaking and recognition have become the most widely accepted of the various hypotheses.
The large size of hadrosaurid eye sockets and the presence of sclerotic rings in the eyes imply acute vision and diurnal habits, evidence that sight was important to these animals. The hadrosaurid sense of hearing also appears to be strong. There is at least one example, in the related Corythosaurus, of a slender stapes (reptilian ear bone) in place, which combined with a large space for an eardrum implies a sensitive middle ear, and the hadrosaurid lagena is elongate like a crocodilian's. This indicates that the auditory portion of the inner ear was well-developed. If used as a noisemaker, the crest could also have provided recognizable differences for different species or sexes, because the differing layouts of the nasal passages corresponding to the different crest shapes would have produced intrinsically different sounds.
Paleoecology
Lambeosaurus lambei and L. magnicristatus, from the Dinosaur Park Formation, were members of a diverse and well-documented fauna of prehistoric animals that included such well-known dinosaurs as the horned Centrosaurus, Styracosaurus, and Chasmosaurus, fellow duckbills Prosaurolophus, Gryposaurus, Corythosaurus, and Parasaurolophus, tyrannosaurid Gorgosaurus, and armored Edmontonia and Euoplocephalus. The Dinosaur Park Formation is interpreted as a low-relief setting of rivers and floodplains that became more swampy and influenced by marine conditions over time as the Western Interior Seaway transgressed westward. The climate was warmer than present-day Alberta, without frost, but with wetter and drier seasons. Conifers were apparently the dominant canopy plants, with an understory of ferns, tree ferns, and angiosperms. The anatomically similar L. lambei, L. magnicristatus, and Corythosaurus were separated by time within the formation, based on stratigraphy. Corythosaurus fossils are known from the lower two-thirds of the Formation, L. lambei fossils are present in the upper third, and L. magnicristatus remains are rare and present only at the very top, where the marine influence was greater.
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