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52135 | https://en.wikipedia.org/wiki/Pneumonia | Pneumonia | Pneumonia is an inflammatory condition of the lung primarily affecting the small air sacs known as alveoli. Symptoms typically include some combination of productive or dry cough, chest pain, fever, and difficulty breathing. The severity of the condition is variable.
Pneumonia is usually caused by infection with viruses or bacteria, and less commonly by other microorganisms. Identifying the responsible pathogen can be difficult. Diagnosis is often based on symptoms and physical examination. Chest X-rays, blood tests, and culture of the sputum may help confirm the diagnosis. The disease may be classified by where it was acquired, such as community- or hospital-acquired or healthcare-associated pneumonia.
Risk factors for pneumonia include cystic fibrosis, chronic obstructive pulmonary disease (COPD), sickle cell disease, asthma, diabetes, heart failure, a history of smoking, a poor ability to cough (such as following a stroke), and immunodeficiency.
Vaccines to prevent certain types of pneumonia (such as those caused by Streptococcus pneumoniae bacteria, influenza viruses, or SARS-CoV-2) are available. Other methods of prevention include hand washing to prevent infection, prompt treatment of worsening respiratory symptoms, and not smoking.
Treatment depends on the underlying cause. Pneumonia believed to be due to bacteria is treated with antibiotics. If the pneumonia is severe, the affected person is generally hospitalized. Oxygen therapy may be used if oxygen levels are low.
Each year, pneumonia affects about 450 million people globally (7% of the population) and results in about 4 million deaths. With the introduction of antibiotics and vaccines in the 20th century, survival has greatly improved. Nevertheless, pneumonia remains a leading cause of death in developing countries, and also among the very old, the very young, and the chronically ill. Pneumonia often shortens the period of suffering among those already close to death and has thus been called "the old man's friend".
Signs and symptoms
People with infectious pneumonia often have a productive cough, fever accompanied by shaking chills, shortness of breath, sharp or stabbing chest pain during deep breaths, and an increased rate of breathing. In elderly people, confusion may be the most prominent sign.
The typical signs and symptoms in children under five are fever, cough, and fast or difficult breathing. Fever is not very specific, as it occurs in many other common illnesses and may be absent in those with severe disease, malnutrition or in the elderly. In addition, a cough is frequently absent in children less than 2 months old. More severe signs and symptoms in children may include blue-tinged skin, unwillingness to drink, convulsions, ongoing vomiting, extremes of temperature, or a decreased level of consciousness.
Bacterial and viral cases of pneumonia usually result in similar symptoms. Some causes are associated with classic, but non-specific, clinical characteristics. Pneumonia caused by Legionella may occur with abdominal pain, diarrhea, or confusion. Pneumonia caused by Streptococcus pneumoniae is associated with rusty colored sputum. Pneumonia caused by Klebsiella may have bloody sputum often described as "currant jelly". Bloody sputum (known as hemoptysis) may also occur with tuberculosis, Gram-negative pneumonia, lung abscesses and more commonly acute bronchitis. Pneumonia caused by Mycoplasma pneumoniae may occur in association with swelling of the lymph nodes in the neck, joint pain, or a middle ear infection. Viral pneumonia presents more commonly with wheezing than bacterial pneumonia. Pneumonia was historically divided into "typical" and "atypical" based on the belief that the presentation predicted the underlying cause. However, evidence has not supported this distinction, therefore it is no longer emphasized.
Cause
Pneumonia is due to infections caused primarily by bacteria or viruses and less commonly by fungi and parasites. Although more than 100 strains of infectious agents have been identified, only a few are responsible for the majority of cases. Mixed infections with both viruses and bacteria may occur in roughly 45% of infections in children and 15% of infections in adults. A causative agent may not be isolated in about half of cases despite careful testing. In an active population-based surveillance for community-acquired pneumonia requiring hospitalization in five hospitals in Chicago and Nashville from January 2010 through June 2012, 2259 patients were identified who had radiographic evidence of pneumonia and specimens that could be tested for the responsible pathogen. Most patients (62%) had no detectable pathogens in their sample, and unexpectedly, respiratory viruses were detected more frequently than bacteria. Specifically, 23% had one or more viruses, 11% had one or more bacteria, 3% had both bacterial and viral pathogens, and 1% had a fungal or mycobacterial infection. "The most common pathogens were human rhinovirus (in 9% of patients), influenza virus (in 6%), and Streptococcus pneumoniae (in 5%)."
The term pneumonia is sometimes more broadly applied to any condition resulting in inflammation of the lungs (caused for example by autoimmune diseases, chemical burns or drug reactions); however, this inflammation is more accurately referred to as pneumonitis.
Factors that predispose to pneumonia include smoking, immunodeficiency, alcoholism, chronic obstructive pulmonary disease, sickle cell disease (SCD), asthma, chronic kidney disease, liver disease, and biological aging. Additional risks in children include not being breastfed, exposure to cigarette smoke and other air pollution, malnutrition, and poverty. The use of acid-suppressing medications – such as proton-pump inhibitors or H2 blockers – is associated with an increased risk of pneumonia. Approximately 10% of people who require mechanical ventilation develop ventilator-associated pneumonia, and people with a gastric feeding tube have an increased risk of developing aspiration pneumonia. Moreover, the misplacement of a feeding tube can lead to aspiration pneumonia. 28% of tube malposition results in pneumonia. As with Avanos Medical's feeding tube placement system, the CORTRAK* 2 EAS, which was recalled in May 2022 by the FDA due to adverse events reported, including pneumonia, caused a total of 60 injuries and 23 patient deaths, as communicated by the FDA. For people with certain variants of the FER gene, the risk of death is reduced in sepsis caused by pneumonia. However, for those with TLR6 variants, the risk of getting Legionnaires' disease is increased.
Bacteria
Bacteria are the most common cause of community-acquired pneumonia (CAP), with Streptococcus pneumoniae isolated in nearly 50% of cases. Other commonly isolated bacteria include Haemophilus influenzae in 20%, Chlamydophila pneumoniae in 13%, and Mycoplasma pneumoniae in 3% of cases; Staphylococcus aureus; Moraxella catarrhalis; and Legionella pneumophila. A number of drug-resistant versions of the above infections are becoming more common, including drug-resistant Streptococcus pneumoniae (DRSP) and methicillin-resistant Staphylococcus aureus (MRSA).
The spreading of organisms is facilitated by certain risk factors. Alcoholism is associated with Streptococcus pneumoniae, anaerobic organisms, and Mycobacterium tuberculosis; smoking facilitates the effects of Streptococcus pneumoniae, Haemophilus influenzae, Moraxella catarrhalis, and Legionella pneumophila. Exposure to birds is associated with Chlamydia psittaci; farm animals with Coxiella burnetti; aspiration of stomach contents with anaerobic organisms; and cystic fibrosis with Pseudomonas aeruginosa and Staphylococcus aureus. Streptococcus pneumoniae is more common in the winter, and it should be suspected in persons aspirating a large number of anaerobic organisms.
Viruses
In adults, viruses account for about one third of pneumonia cases, and in children for about 15% of them. Commonly implicated agents include rhinoviruses, coronaviruses, influenza virus, respiratory syncytial virus (RSV), adenovirus, and parainfluenza. Herpes simplex virus rarely causes pneumonia, except in groups such as newborns, persons with cancer, transplant recipients, and people with significant burns. After organ transplantation or in otherwise immunocompromised persons, there are high rates of cytomegalovirus pneumonia. Those with viral infections may be secondarily infected with the bacteria Streptococcus pneumoniae, Staphylococcus aureus, or Haemophilus influenzae, particularly when other health problems are present. Different viruses predominate at different times of the year; during flu season, for example, influenza may account for more than half of all viral cases. Outbreaks of other viruses also occur occasionally, including hantaviruses and coronaviruses. Severe acute respiratory syndrome coronavirus 2 (SARS-CoV-2) can also result in pneumonia.
Fungi
Fungal pneumonia is uncommon, but occurs more commonly in individuals with weakened immune systems due to AIDS, immunosuppressive drugs, or other medical problems. It is most often caused by Histoplasma capsulatum, Blastomyces, Cryptococcus neoformans, Pneumocystis jiroveci (pneumocystis pneumonia, or PCP), and Coccidioides immitis. Histoplasmosis is most common in the Mississippi River basin, and coccidioidomycosis is most common in the Southwestern United States. The number of cases of fungal pneumonia has been increasing in the latter half of the 20th century due to increasing travel and rates of immunosuppression in the population. For people infected with HIV/AIDS, PCP is a common opportunistic infection.
Parasites
A variety of parasites can affect the lungs, including Toxoplasma gondii, Strongyloides stercoralis, Ascaris lumbricoides, and Plasmodium malariae. These organisms typically enter the body through direct contact with the skin, ingestion, or via an insect vector. Except for Paragonimus westermani, most parasites do not specifically affect the lungs but involve the lungs secondarily to other sites. Some parasites, in particular those belonging to the Ascaris and Strongyloides genera, stimulate a strong eosinophilic reaction, which may result in eosinophilic pneumonia. In other infections, such as malaria, lung involvement is due primarily to cytokine-induced systemic inflammation. In the developed world, these infections are most common in people returning from travel or in immigrants. Around the world, parasitic pneumonia is most common in the immunodeficient.
Noninfectious
Idiopathic interstitial pneumonia or noninfectious pneumonia is a class of diffuse lung diseases. They include diffuse alveolar damage, organizing pneumonia, nonspecific interstitial pneumonia, lymphocytic interstitial pneumonia, desquamative interstitial pneumonia, respiratory bronchiolitis interstitial lung disease, and usual interstitial pneumonia. Lipoid pneumonia is another rare cause due to lipids entering the lung. These lipids can either be inhaled or spread to the lungs from elsewhere in the body.
Mechanisms
Pneumonia frequently starts as an upper respiratory tract infection that moves into the lower respiratory tract. It is a type of pneumonitis (lung inflammation). The normal flora of the upper airway give protection by competing with pathogens for nutrients. In the lower airways, reflexes of the glottis, actions of complement proteins and immunoglobulins are important for protection. Microaspiration of contaminated secretions can infect the lower airways and cause pneumonia. The progress of pneumonia is determined by the virulence of the organism; the amount of organism required to start an infection; and the body's immune response against the infection.
Bacterial
Most bacteria enter the lungs via small aspirations of organisms residing in the throat or nose. Half of normal people have these small aspirations during sleep. While the throat always contains bacteria, potentially infectious ones reside there only at certain times and under certain conditions. A minority of types of bacteria such as Mycobacterium tuberculosis and Legionella pneumophila reach the lungs via contaminated airborne droplets. Bacteria can also spread via the blood. Once in the lungs, bacteria may invade the spaces between cells and between alveoli, where the macrophages and neutrophils (defensive white blood cells) attempt to inactivate the bacteria. The neutrophils also release cytokines, causing a general activation of the immune system. This leads to the fever, chills, and fatigue common in bacterial pneumonia. The neutrophils, bacteria, and fluid from surrounding blood vessels fill the alveoli, resulting in the consolidation seen on chest X-ray.
Viral
Viruses may reach the lung by a number of different routes. Respiratory syncytial virus is typically contracted when people touch contaminated objects and then touch their eyes or nose. Other viral infections occur when contaminated airborne droplets are inhaled through the nose or mouth. Once in the upper airway, the viruses may make their way into the lungs, where they invade the cells lining the airways, alveoli, or lung parenchyma. Some viruses such as measles and herpes simplex may reach the lungs via the blood. The invasion of the lungs may lead to varying degrees of cell death. When the immune system responds to the infection, even more lung damage may occur. Primarily white blood cells, mainly mononuclear cells, generate the inflammation. As well as damaging the lungs, many viruses simultaneously affect other organs and thus disrupt other body functions. Viruses also make the body more susceptible to bacterial infections; in this way, bacterial pneumonia can occur at the same time as viral pneumonia.
Diagnosis
Pneumonia is typically diagnosed based on a combination of physical signs and often a chest X-ray. In recent years, however, the role of lung ultrasonography has gained prominence, with substantial evidence demonstrating that, in expert hands, it surpasses radiography in accuracy. In adults with normal vital signs and a normal lung examination, the diagnosis is unlikely. However, the underlying cause can be difficult to confirm, as there is no definitive test able to distinguish between bacterial and non-bacterial cause. The overall impression of a physician appears to be at least as good as decision rules for making or excluding the diagnosis.
Diagnosis in children
The World Health Organization has defined pneumonia in children clinically based on either a cough or difficulty breathing and a rapid respiratory rate, chest indrawing, or a decreased level of consciousness. A rapid respiratory rate is defined as greater than 60 breaths per minute in children under 2 months old, greater than 50 breaths per minute in children 2 months to 1 year old, or greater than 40 breaths per minute in children 1 to 5 years old.
In children, low oxygen levels and lower chest indrawing are more sensitive than hearing chest crackles with a stethoscope or increased respiratory rate. Grunting and nasal flaring may be other useful signs in children less than five years old.
Lack of wheezing is an indicator of Mycoplasma pneumoniae in children with pneumonia, but as an indicator it is not accurate enough to decide whether or not macrolide treatment should be used. The presence of chest pain in children with pneumonia doubles the probability of Mycoplasma pneumoniae.
Diagnosis in adults
In general, in adults, investigations are not needed in mild cases. There is a very low risk of pneumonia if all vital signs and auscultation are normal. C-reactive protein (CRP) may help support the diagnosis. For those with CRP less than 20 mg/L without convincing evidence of pneumonia, antibiotics are not recommended.
Procalcitonin may help determine the cause and support decisions about who should receive antibiotics. Antibiotics are encouraged if the procalcitonin level reaches 0.25 μg/L, strongly encouraged if it reaches 0.5 μg/L, and strongly discouraged if the level is below 0.10 μg/L. In people requiring hospitalization, pulse oximetry, chest radiography and blood tests – including a complete blood count, serum electrolytes, C-reactive protein level, and possibly liver function tests – are recommended.
The diagnosis of influenza-like illness can be made based on the signs and symptoms; however, confirmation of an influenza infection requires testing. Thus, treatment is frequently based on the presence of influenza in the community or a rapid influenza test.
Adults 65 years old or older, as well as cigarette smokers and people with ongoing medical conditions are at increased risk for pneumonia.
Physical exam
Physical examination may sometimes reveal low blood pressure, high heart rate, or low oxygen saturation. The respiratory rate may be faster than normal, and this may occur a day or two before other signs. Examination of the chest may be normal, but it may show decreased expansion on the affected side. Harsh breath sounds from the larger airways that are transmitted through the inflamed lung are termed bronchial breathing and are heard on auscultation with a stethoscope. Crackles (rales) may be heard over the affected area during inspiration. Percussion may be dulled over the affected lung, and increased, rather than decreased, vocal resonance distinguishes pneumonia from a pleural effusion.
Imaging
A chest radiograph is frequently used in diagnosis. In people with mild disease, imaging is needed only in those with potential complications, those not having improved with treatment, or those in which the cause is uncertain. If a person is sufficiently sick to require hospitalization, a chest radiograph is recommended. Findings do not always match the severity of disease and do not reliably separate between bacterial and viral infection.
X-ray presentations of pneumonia may be classified as lobar pneumonia, bronchopneumonia, lobular pneumonia, and interstitial pneumonia. Bacterial, community-acquired pneumonia classically show lung consolidation of one lung segmental lobe, which is known as lobar pneumonia. However, findings may vary, and other patterns are common in other types of pneumonia. Aspiration pneumonia may present with bilateral opacities primarily in the bases of the lungs and on the right side. Radiographs of viral pneumonia may appear normal, appear hyper-inflated, have bilateral patchy areas, or present similar to bacterial pneumonia with lobar consolidation. Radiologic findings may not be present in the early stages of the disease, especially in the presence of dehydration, or may be difficult to interpret in the obese or those with a history of lung disease. Complications such as pleural effusion may also be found on chest radiographs. Laterolateral chest radiographs can increase the diagnostic accuracy of lung consolidation and pleural effusion.
A CT scan can give additional information in indeterminate cases and provide more details in those with an unclear chest radiograph (for example occult pneumonia in chronic obstructive pulmonary disease). They can be used to exclude pulmonary embolism and fungal pneumonia, and detect lung abscesses in those who are not responding to treatments. However, CT scans are more expensive, have a higher dose of radiation, and cannot be done at bedside.
Lung ultrasound may also be useful in helping to make the diagnosis. Ultrasound is radiation free and can be done at bedside. However, ultrasound requires specific skills to operate the machine and interpret the findings. It may be more accurate than chest X-ray.
Microbiology
In people managed in the community, determining the causative agent is not cost-effective and typically does not alter management. For people who do not respond to treatment, sputum culture should be considered, and culture for Mycobacterium tuberculosis should be carried out in persons with a chronic productive cough. Microbiological evaluation is also indicated in severe pneumonia, alcoholism, asplenia, immunosuppression, HIV infection, and those being empirically treated for MRSA of pseudomonas. Although positive blood culture and pleural fluid culture definitively establish the diagnosis of the type of micro-organism involved, a positive sputum culture has to be interpreted with care for the possibility of colonisation of respiratory tract. Testing for other specific organisms may be recommended during outbreaks, for public health reasons. In those hospitalized for severe disease, both sputum and blood cultures are recommended, as well as testing the urine for antigens to Legionella and Streptococcus. Viral infections, can be confirmed via detection of either the virus or its antigens with culture or polymerase chain reaction (PCR), among other techniques. Mycoplasma, Legionella, Streptococcus, and Chlamydia can also be detected using PCR techniques on bronchoalveolar lavage and nasopharyngeal swab. The causative agent is determined in only 15% of cases with routine microbiological tests.
Classification
Pneumonitis refers to lung inflammation; pneumonia refers to pneumonitis, usually due to infection but sometimes non-infectious, that has the additional feature of pulmonary consolidation. Pneumonia is most commonly classified by where or how it was acquired: community-acquired, aspiration, healthcare-associated, hospital-acquired, and ventilator-associated pneumonia. It may also be classified by the area of the lung affected: lobar, bronchial pneumonia and acute interstitial pneumonia; or by the causative organism. Pneumonia in children may additionally be classified based on signs and symptoms as non-severe, severe, or very severe.
The setting in which pneumonia develops is important to treatment, as it correlates to which pathogens are likely suspects, which mechanisms are likely, which antibiotics are likely to work or fail, and which complications can be expected based on the person's health status.
Community
Community-acquired pneumonia (CAP) is acquired in the community, outside of health care facilities. Compared with healthcare-associated pneumonia, it is less likely to involve multidrug-resistant bacteria. Although the latter are no longer rare in CAP, they are still less likely. Prior stays in healthcare-related environments such as hospitals, nursing homes, or hemodialysis centers or a history of receiving domiciliary care can increase patients' risk for CAP caused by multidrug-resistant bacteria.
Healthcare
Health care–associated pneumonia (HCAP) is an infection associated with recent exposure to the health care system, including hospitals, outpatient clinics, nursing homes, dialysis centers, chemotherapy treatment, or home care. HCAP is sometimes called MCAP (medical care–associated pneumonia).
People may become infected with pneumonia in a hospital; this is defined as pneumonia not present at the time of admission (symptoms must start at least 48 hours after admission). It is likely to involve hospital-acquired infections, with higher risk of multidrug-resistant pathogens. People in a hospital often have other medical conditions, which may make them more susceptible to pathogens in the hospital.
Ventilator-associated pneumonia occurs in people breathing with the help of mechanical ventilation. Ventilator-associated pneumonia is specifically defined as pneumonia that arises more than 48 to 72 hours after endotracheal intubation.
Differential diagnosis
Several diseases can present with similar signs and symptoms to pneumonia, such as: chronic obstructive pulmonary disease, asthma, pulmonary edema, bronchiectasis, lung cancer, and pulmonary emboli. Unlike pneumonia, asthma and COPD typically present with wheezing, pulmonary edema presents with an abnormal electrocardiogram, cancer and bronchiectasis present with a cough of longer duration, and pulmonary emboli present with acute onset sharp chest pain and shortness of breath. Mild pneumonia should be differentiated from upper respiratory tract infection (URTI). Severe pneumonia should be differentiated from acute heart failure. Pulmonary infiltrates that resolved after giving mechanical ventilation should point to heart failure and atelectasis rather than pneumonia. For recurrent pneumonia, underlying lung cancer, metastasis, tuberculosis, a foreign bodies, immunosuppression, and hypersensitivity should be suspected.
Prevention
Prevention includes vaccination, environmental measures, and appropriate treatment of other health problems. It is believed that, if appropriate preventive measures were instituted globally, mortality among children could be reduced by 400,000; and, if proper treatment were universally available, childhood deaths could be decreased by another 600,000.
Vaccination
Vaccination prevents against certain bacterial and viral pneumonias both in children and adults. Influenza vaccines are modestly effective at preventing symptoms of influenza, The Centers for Disease Control and Prevention (CDC) recommends yearly influenza vaccination for every person 6 months and older. Immunizing health care workers decreases the risk of viral pneumonia among their patients.
Vaccinations against Haemophilus influenzae and Streptococcus pneumoniae have good evidence to support their use. There is strong evidence for vaccinating children under the age of 2 against Streptococcus pneumoniae (pneumococcal conjugate vaccine). Vaccinating children against Streptococcus pneumoniae has led to a decreased rate of these infections in adults, because many adults acquire infections from children. A Streptococcus pneumoniae vaccine is available for adults, and has been found to decrease the risk of invasive pneumococcal disease by 74%, but there is insufficient evidence to suggest using the pneumococcal vaccine to prevent pneumonia or death in the general adult population. The CDC recommends that young children and adults over the age of 65 receive the pneumococcal vaccine, as well as older children or younger adults who have an increased risk of getting pneumococcal disease. The pneumococcal vaccine has been shown to reduce the risk of community acquired pneumonia in people with chronic obstructive pulmonary disease, but does not reduce mortality or the risk of hospitalization for people with this condition. People with COPD are recommended by a number of guidelines to have a pneumococcal vaccination. Other vaccines for which there is support for a protective effect against pneumonia include pertussis, varicella, and measles.
Medications
When influenza outbreaks occur, medications such as amantadine or rimantadine may help prevent the condition, but they are associated with side effects. Zanamivir or oseltamivir decrease the chance that people who are exposed to the virus will develop symptoms; however, it is recommended that potential side effects are taken into account.
Other
Smoking cessation and reducing indoor air pollution, such as that from cooking indoors with wood, crop residues or dung, are both recommended. Smoking appears to be the single biggest risk factor for pneumococcal pneumonia in otherwise-healthy adults. Hand hygiene and coughing into one's sleeve may also be effective preventative measures. Wearing surgical masks by the sick may also prevent illness.
Appropriately treating underlying illnesses (such as HIV/AIDS, diabetes mellitus, and malnutrition) can decrease the risk of pneumonia. In children less than 6 months of age, exclusive breast feeding reduces both the risk and severity of disease. In people with HIV/AIDS and a CD4 count of less than 200 cells/uL the antibiotic trimethoprim/sulfamethoxazole decreases the risk of Pneumocystis pneumonia and is also useful for prevention in those that are immunocompromised but do not have HIV.
Testing pregnant women for Group B Streptococcus and Chlamydia trachomatis, and administering antibiotic treatment, if needed, reduces rates of pneumonia in infants; preventive measures for HIV transmission from mother to child may also be efficient. Suctioning the mouth and throat of infants with meconium-stained amniotic fluid has not been found to reduce the rate of aspiration pneumonia and may cause potential harm, thus this practice is not recommended in the majority of situations. In the frail elderly good oral health care may lower the risk of aspiration pneumonia, even though there is no good evidence that one approach to mouth care is better than others in preventing nursing home acquired pneumonia. Zinc supplementation in children 2 months to five years old appears to reduce rates of pneumonia.
For people with low levels of vitamin C in their diet or blood, taking vitamin C supplements may be suggested to decrease the risk of pneumonia, although there is no strong evidence of benefit. There is insufficient evidence to recommend that the general population take vitamin C to prevent or treat pneumonia.
For adults and children in the hospital who require a respirator, there is no strong evidence indicating a difference between heat and moisture exchangers and heated humidifiers for preventing pneumonia. There is tentative evidence that laying flat on the back compared to semi-raised increases pneumonia risks in people who are intubated.
Management
Antibiotics by mouth, rest, simple analgesics, and fluids usually suffice for complete resolution. However, those with other medical conditions, the elderly, or those with significant trouble breathing may require more advanced care. If the symptoms worsen, the pneumonia does not improve with home treatment, or complications occur, hospitalization may be required. Worldwide, approximately 7–13% of cases in children result in hospitalization, whereas in the developed world between 22 and 42% of adults with community-acquired pneumonia are admitted. The CURB-65 score is useful for determining the need for admission in adults. If the score is 0 or 1, people can typically be managed at home; if it is 2, a short hospital stay or close follow-up is needed; if it is 3–5, hospitalization is recommended. In children those with respiratory distress or oxygen saturations of less than 90% should be hospitalized. The utility of chest physiotherapy in pneumonia has not yet been determined. Over-the-counter cough medicine has not been found to be effective, nor has the use of zinc supplementation in children. There is insufficient evidence for mucolytics. There is no strong evidence to recommend that children who have non-measles related pneumonia take vitamin A supplements. Vitamin D, as of 2023, is of unclear benefit in children. Vitamin C administration in pneumonia needs further research, although it can be given to patient of low plasma vitamin C because it is not expensive and low risk.
Pneumonia can cause severe illness in a number of ways, and pneumonia with evidence of organ dysfunction may require intensive care unit admission for observation and specific treatment. The main impact is on the respiratory and the circulatory system. Respiratory failure not responding to normal oxygen therapy may require heated humidified high-flow therapy delivered through nasal cannulae, non-invasive ventilation, or in severe cases mechanical ventilation through an endotracheal tube. Regarding circulatory problems as part of sepsis, evidence of poor blood flow or low blood pressure is initially treated with 30 mL/kg of crystalloid infused intravenously. In situations where fluids alone are ineffective, vasopressor medication may be required.
For adults with moderate or severe acute respiratory distress syndrome (ARDS) undergoing mechanical ventilation, there is a reduction in mortality when people lie on their front for at least 12 hours a day. However, this increases the risk of endotracheal tube obstruction and pressure sores.
Bacterial
Antibiotics improve outcomes in those with bacterial pneumonia. The first dose of antibiotics should be given as soon as possible. Increased use of antibiotics, however, may lead to the development of antimicrobial resistant strains of bacteria. Antibiotic choice depends initially on the characteristics of the person affected, such as age, underlying health, and the location the infection was acquired. Antibiotic use is also associated with side effects such as nausea, diarrhea, dizziness, taste distortion, or headaches. In the UK, treatment before culture results with amoxicillin is recommended as the first line for community-acquired pneumonia, with doxycycline or clarithromycin as alternatives. In North America, amoxicillin, doxycycline, and in some areas a macrolide (such as azithromycin or erythromycin) is the first-line outpatient treatment in adults. In children with mild or moderate symptoms, amoxicillin taken by mouth is the first line. The use of fluoroquinolones in uncomplicated cases is discouraged due to concerns about side-effects and generating resistance in light of there being no greater benefit.
For those who require hospitalization and caught their pneumonia in the community the use of a β-lactam such as cephazolin plus a macrolide such as azithromycin is recommended. A fluoroquinolone may replace azithromycin but is less preferred. Antibiotics by mouth and by injection appear to be similarly effective in children with severe pneumonia.
The duration of treatment has traditionally been seven to ten days, but increasing evidence suggests that shorter courses (3–5 days) may be effective for certain types of pneumonia and may reduce the risk of antibiotic resistance. Research in children showed that a shorter, 3-day course of amoxicillin was as effective as a longer, 7-day course for treating pneumonia in this population. For pneumonia that is associated with a ventilator caused by non-fermenting Gram-negative bacilli (NF-GNB), a shorter course of antibiotics increases the risk that the pneumonia will return. Recommendations for hospital-acquired pneumonia include third- and fourth-generation cephalosporins, carbapenems, fluoroquinolones, aminoglycosides, and vancomycin. These antibiotics are often given intravenously and used in combination. In those treated in hospital, more than 90% improve with the initial antibiotics. For people with ventilator-acquired pneumonia, the choice of antibiotic therapy will depend on the person's risk of being infected with a strain of bacteria that is multi-drug resistant. Once clinically stable, intravenous antibiotics should be switched to oral antibiotics. For those with Methicillin resistant Staphylococcus aureus (MRSA) or Legionella infections, prolonged antibiotics may be beneficial.
The addition of corticosteroids to standard antibiotic treatment appears to improve outcomes, reducing death and morbidity for adults with severe community acquired pneumonia, and reducing death for adults and children with non-severe community acquired pneumonia. A 2017 review therefore recommended them in adults with severe community acquired pneumonia. A 2019 guideline however recommended against their general use, unless refractory shock was present. Side effects associated with the use of corticosteroids include high blood sugar. There is some evidence that adding corticosteroids to the standard PCP pneumonia treatment may be beneficial for people who are infected with HIV.
The use of granulocyte colony stimulating factor (G-CSF) along with antibiotics does not appear to reduce mortality and routine use for treating pneumonia is not supported by evidence.
Viral
Neuraminidase inhibitors may be used to treat viral pneumonia caused by influenza viruses (influenza A and influenza B). No specific antiviral medications are recommended for other types of community acquired viral pneumonias including SARS coronavirus, adenovirus, hantavirus, and parainfluenza virus. Influenza A may be treated with rimantadine or amantadine, while influenza A or B may be treated with oseltamivir, zanamivir or peramivir. These are of most benefit if they are started within 48 hours of the onset of symptoms. Many strains of H5N1 influenza A, also known as avian influenza or "bird flu", have shown resistance to rimantadine and amantadine. The use of antibiotics in viral pneumonia is recommended by some experts, as it is impossible to rule out a complicating bacterial infection. The British Thoracic Society recommends that antibiotics be withheld in those with mild disease. The use of corticosteroids is controversial.
Aspiration
In general, aspiration pneumonitis is treated conservatively with antibiotics indicated only for aspiration pneumonia. The choice of antibiotic will depend on several factors, including the suspected causative organism and whether pneumonia was acquired in the community or developed in a hospital setting. Common options include clindamycin, a combination of a beta-lactam antibiotic and metronidazole, or an aminoglycoside.
Corticosteroids are sometimes used in aspiration pneumonia, but there is limited evidence to support their effectiveness.
Follow-up
The British Thoracic Society recommends that a follow-up chest radiograph be taken in people with persistent symptoms, smokers, and people older than 50. American guidelines vary, from generally recommending a follow-up chest radiograph to not mentioning any follow-up.
Prognosis
With treatment, most types of bacterial pneumonia will stabilize in 3–6 days. It often takes a few weeks before most symptoms resolve. X-ray findings typically clear within four weeks and mortality is low (less than 1%). In the elderly or people with other lung problems, recovery may take more than 12 weeks. In persons requiring hospitalization, mortality may be as high as 10%, and in those requiring intensive care it may reach 30–50%. Pneumonia is the most common hospital-acquired infection that causes death. Before the advent of antibiotics, mortality was typically 30% in those that were hospitalized. However, for those whose lung condition deteriorates within 72 hours, the problem is usually due to sepsis. If pneumonia deteriorates after 72 hours, it could be due to nosocomial infection or excerbation of other underlying comorbidities. About 10% of those discharged from hospital are readmitted due to underlying co-morbidities such as heart, lung, or neurological disorders, or due to new onset of pneumonia.
Complications may occur in particular in the elderly and those with underlying health problems. This may include, among others: empyema, lung abscess, bronchiolitis obliterans, acute respiratory distress syndrome, sepsis, and worsening of underlying health problems.
Clinical prediction rules
Clinical prediction rules have been developed to more objectively predict outcomes of pneumonia. These rules are often used to decide whether to hospitalize the person.
CURB-65 score, which takes into account the severity of symptoms, any underlying diseases, and age
Pneumonia severity index (or PSI Score)
Pleural effusion, empyema, and abscess
In pneumonia, a collection of fluid may form in the space that surrounds the lung. Occasionally, microorganisms will infect this fluid, causing an empyema. To distinguish an empyema from the more common simple parapneumonic effusion, the fluid may be collected with a needle (thoracentesis), and examined. If this shows evidence of empyema, complete drainage of the fluid is necessary, often requiring a drainage catheter. In severe cases of empyema, surgery may be needed. If the infected fluid is not drained, the infection may persist, because antibiotics do not penetrate well into the pleural cavity. If the fluid is sterile, it must be drained only if it is causing symptoms or remains unresolved.
In rare circumstances, bacteria in the lung will form a pocket of infected fluid called a lung abscess. Lung abscesses can usually be seen with a chest X-ray but frequently require a chest CT scan to confirm the diagnosis. Abscesses typically occur in aspiration pneumonia, and often contain several types of bacteria. Long-term antibiotics are usually adequate to treat a lung abscess, but sometimes the abscess must be drained by a surgeon or radiologist.
Respiratory and circulatory failure
Pneumonia can cause respiratory failure by triggering acute respiratory distress syndrome (ARDS), which results from a combination of infection and inflammatory response. The lungs quickly fill with fluid and become stiff. This stiffness, combined with severe difficulties extracting oxygen due to the alveolar fluid, may require long periods of mechanical ventilation for survival. Other causes of circulatory failure are hypoxemia, inflammation, and increased coagulability.
Sepsis is a potential complication of pneumonia but usually occurs in people with poor immunity or hyposplenism. The organisms most commonly involved are Streptococcus pneumoniae, Haemophilus influenzae, and Klebsiella pneumoniae. Other causes of the symptoms should be considered such as a myocardial infarction or a pulmonary embolism.
Epidemiology
Pneumonia is a common illness affecting approximately 450 million people a year and occurring in all parts of the world. It is a major cause of death among all age groups resulting in 4 million deaths (7% of the world's total death) yearly. Rates are greatest in children less than five, and adults older than 75 years. It occurs about five times more frequently in the developing world than in the developed world. Viral pneumonia accounts for about 200 million cases. In the United States, , pneumonia is the 8th leading cause of death.
Children
In 2008, pneumonia occurred in approximately 156 million children (151 million in the developing world and 5 million in the developed world). In 2010, it resulted in 1.3 million deaths, or 18% of all deaths in those under five years, of which 95% occurred in the developing world. Countries with the greatest burden of disease include India (43 million), China (21 million) and Pakistan (10 million). It is the leading cause of death among children in low income countries. Many of these deaths occur in the newborn period. The World Health Organization estimates that one in three newborn infant deaths is due to pneumonia. Approximately half of these deaths can be prevented, as they are caused by the bacteria for which an effective vaccine is available. The IDSA has recommended that children and infants with symptoms of CAP should be hospitalized so they have access to pediatric nursing care. In 2011, pneumonia was the most common reason for admission to the hospital after an emergency department visit in the U.S. for infants and children.
History
Pneumonia has been a common disease throughout human history. The word is from Greek πνεύμων (pneúmōn) meaning "lung". The symptoms were described by Hippocrates (–370 BC): "Peripneumonia, and pleuritic affections, are to be thus observed: If the fever be acute, and if there be pains on either side, or in both, and if expiration be if cough be present, and the sputa expectorated be of a blond or livid color, or likewise thin, frothy, and florid, or having any other character different from the common... When pneumonia is at its height, the case is beyond remedy if he is not purged, and it is bad if he has dyspnoea, and urine that is thin and acrid, and if sweats come out about the neck and head, for such sweats are bad, as proceeding from the suffocation, rales, and the violence of the disease which is obtaining the upper hand." However, Hippocrates referred to pneumonia as a disease "named by the ancients". He also reported the results of surgical drainage of empyemas. Maimonides (1135–1204 AD) observed: "The basic symptoms that occur in pneumonia and that are never lacking are as follows: acute fever, sticking pleuritic pain in the side, short rapid breaths, serrated pulse and cough." This clinical description is quite similar to those found in modern textbooks, and it reflected the extent of medical knowledge through the Middle Ages into the 19th century.
Edwin Klebs was the first to observe bacteria in the airways of persons having died of pneumonia in 1875. Initial work identifying the two common bacterial causes, Streptococcus pneumoniae and Klebsiella pneumoniae, was performed by Carl Friedländer and Albert Fraenkel in 1882 and 1884, respectively. Friedländer's initial work introduced the Gram stain, a fundamental laboratory test still used today to identify and categorize bacteria. Christian Gram's paper describing the procedure in 1884 helped to differentiate the two bacteria, and showed that pneumonia could be caused by more than one microorganism. In 1887, Jaccond demonstrated pneumonia may be caused by opportunistic bacteria always present in the lung.
Sir William Osler, known as "the father of modern medicine", appreciated the death and disability caused by pneumonia, describing it as the "captain of the men of death" in 1918, as it had overtaken tuberculosis as one of the leading causes of death at the time. This phrase was originally coined by John Bunyan in reference to "consumption" (tuberculosis). Osler also described pneumonia as "the old man's friend" as death was often quick and painless when there were much slower and more painful ways to die.
Viral pneumonia was first described by Hobart Reimann in 1938. Reimann, Chairman of the Department of Medicine at Jefferson Medical College, had established the practice of routinely typing the pneumococcal organism in cases where pneumonia presented. Out of this work, the distinction between viral and bacterial strains was noticed.
Several developments in the 1900s improved the outcome for those with pneumonia. With the advent of penicillin and other antibiotics, modern surgical techniques, and intensive care in the 20th century, mortality from pneumonia, which had approached 30%, dropped precipitously in the developed world. Vaccination of infants against Haemophilus influenzae type B began in 1988 and led to a dramatic decline in cases shortly thereafter. Vaccination against Streptococcus pneumoniae in adults began in 1977, and in children in 2000, resulting in a similar decline.
Society and culture
Awareness
Due to the relatively low awareness of the disease, 12 November was declared in 2009 as the annual World Pneumonia Day, a day for concerned citizens and policy makers to take action against the disease.
Costs
The global economic cost of community-acquired pneumonia has been estimated at $17 billion annually. Other estimates are considerably higher. In 2012 the estimated aggregate costs of treating pneumonia in the United States were $20 billion; the median cost of a single pneumonia-related hospitalization is over $15,000. According to data released by the Centers for Medicare and Medicaid Services, average 2012 hospital charges for inpatient treatment of uncomplicated pneumonia in the U.S. were $24,549 and ranged as high as $124,000. The average cost of an emergency room consult for pneumonia was $943 and the average cost for medication was $66. Aggregate annual costs of treating pneumonia in Europe have been estimated at €10 billion.
| Biology and health sciences | Illness and injury | null |
52136 | https://en.wikipedia.org/wiki/Citrus | Citrus | Citrus is a genus of flowering trees and shrubs in the family Rutaceae. Plants in the genus produce citrus fruits, including important crops such as oranges, mandarins, lemons, grapefruits, pomelos, and limes.
Citrus is native to South Asia, East Asia, Southeast Asia, Melanesia, and Australia. Indigenous people in these areas have used and domesticated various species since ancient times. Its cultivation first spread into Micronesia and Polynesia through the Austronesian expansion (–1500 BCE). Later, it was spread to the Middle East and the Mediterranean () via the incense trade route, and from Europe to the Americas.
Renowned for their highly fragrant aromas and complex flavor, citrus are among the most popular fruits in cultivation. With a propensity to hybridize between species, making their taxonomy complicated, there are numerous varieties encompassing a wide range of appearance and fruit flavors.
Evolution
Evolutionary history
The large citrus fruit of today evolved originally from small, edible berries over millions of years. Citrus species began to diverge from a common ancestor about 15 million years ago, at about the same time that Severinia (such as the Chinese box orange) diverged from the same ancestor. About 7 million years ago, the ancestors of Citrus split into the main genus, Citrus, and the Poncirus group (such as the trifoliate orange), which some taxonomies consider a separate genus and others include in Citrus Poncirus
is closely enough related that it can still be hybridized with all other citrus and used as rootstock. These estimates are made using genetic mapping of plant chloroplasts. A DNA study published in Nature in 2018 concludes that the genus Citrus evolved in the foothills of the Himalayas, in the area of Assam (India), western Yunnan (China), and northern Myanmar.
The three ancestral species in the genus Citrus associated with modern Citrus cultivars are the mandarin orange, pomelo, and citron. Almost all of the common commercially important citrus fruits (sweet oranges, lemons, grapefruit, limes, and so on) are hybrids between these three species, their main progenies, and other wild Citrus species within the last few thousand years.
Citrus plants are native to subtropical and tropical regions of Asia, Island Southeast Asia, Near Oceania, and northeastern and central Australia. Domestication of citrus species involved much hybridization and introgression, leaving much uncertainty about when and where domestication first happened. A genomic, phylogenic, and biogeographical analysis by Wu et al. (2018) has shown that the center of origin of the genus Citrus is likely the southeast foothills of the Himalayas, in a region stretching from eastern Assam, northern Myanmar, to western Yunnan. It diverged from a common ancestor with Poncirus trifoliata. A change in climate conditions during the Late Miocene (11.63 to 5.33 mya) resulted in a sudden speciation event. The species resulting from this event include the citrons (Citrus medica) of South Asia; the pomelos (C. maxima) of Mainland Southeast Asia; the mandarins (C. reticulata), kumquats (C. japonica), mangshanyegan (C. mangshanensis), and ichang papedas (C. cavaleriei) of southeastern China; the kaffir limes (C. hystrix) of Island Southeast Asia; and the biasong and samuyao (C. micrantha) of the Philippines.
This was followed by the spread of citrus species into Taiwan and Japan in the Early Pliocene (5.33 to 3.6 mya), resulting in the tachibana orange (C. tachibana); and beyond the Wallace Line into Papua New Guinea and Australia during the Early Pleistocene (2.5 million to 800,000 years ago), where further speciation events created in the Australian limes.
Fossil record
A fossil leaf from the Pliocene of Valdarno, Italy is described as †Citrus meletensis.
In China, fossil leaf specimens of †Citrus linczangensis have been collected from late Miocene coal-bearing strata of the Bangmai Formation in Yunnan province. C. linczangensis resembles C. meletensis in having an intramarginal vein, an entire margin, and an articulated and distinctly winged petiole.
Taxonomy
Many cultivated Citrus species are natural or artificial hybrids of a small number of core ancestral species, including the citron, pomelo, and mandarin. Natural and cultivated citrus hybrids include commercially important fruit such as oranges, grapefruit, lemons, limes, and some tangerines. The multiple hybridisations have made the taxonomy of Citrus complex.
Apart from these core species, Australian limes and the recently discovered mangshanyegan are grown. Kumquats and Clymenia spp. are now generally considered to belong within the genus Citrus. The false oranges, Oxanthera from New Caledonia, have been transferred to the Citrus genus on phylogenetic evidence. A recent taxonomy reincorporates the trifoliate orange (Poncirus) into an enlarged Citrus, but recognizes that many botanists still follow Swingle in splitting it off.
History
The earliest introductions of citrus species by human migrations was during the Austronesian expansion (–1500 BCE), where Citrus hystrix, Citrus macroptera, and Citrus maxima were among the canoe plants carried by Austronesian voyagers eastwards into Micronesia and Polynesia.
The citron (Citrus medica) was also introduced early into the Mediterranean basin from India and Southeast Asia. It was introduced via two ancient trade routes: an overland route through Persia, the Levant and the Mediterranean islands; and a maritime route through the Arabian Peninsula and Ptolemaic Egypt into North Africa. Although the exact date of the original introduction is unknown due to the sparseness of archaeobotanical remains, the earliest evidence are seeds recovered from the Hala Sultan Tekke site of Cyprus, dated to around 1200 BCE. Other archaeobotanical evidence includes pollen from Carthage dating back to the 4th century BCE; and carbonized seeds from Pompeii dated to around the 3rd to 2nd century BCE. The earliest complete description of the citron was written by Theophrastus, .
Lemons, pomelos, and sour oranges were introduced to the Mediterranean by Arab traders around the 10th century CE. Sweet oranges were brought to Europe by the Genoese and Portuguese from Asia during the 15th to 16th century. Mandarins were not introduced until the 19th century. Oranges were introduced to Florida by Spanish colonists. In cooler parts of Europe, citrus fruit was grown in orangeries starting in the 17th century; many were as much status symbols as functional agricultural structures.
Etymology
The generic name Citrus originates from Latin, where it denoted either the citron (C. medica) or a conifer tree (Thuja). The Latin word is related to the ancient Greek word for the cedar of Lebanon, (), perhaps from a perceived similarity of the smell of citrus leaves and fruit with that of cedar.
Description
Tree
Citrus plants are large shrubs or small to moderate-sized trees, reaching tall, with spiny shoots and alternately arranged evergreen leaves with an entire margin. The flowers are solitary or in small corymbs, each flower diameter, with five (rarely four) white petals and numerous stamens; they are often very strongly scented, due to the presence of essential oil glands.
Fruit
The fruit is a hesperidium, a specialised berry with multiple carpels, globose to elongated, long and diameter, with a leathery rind or "peel" called a pericarp. The outermost layer of the pericarp is an "exocarp" called the flavedo, commonly referred to as the zest. The middle layer of the pericarp is the mesocarp, which in citrus fruits consists of the white, spongy albedo or pith. The innermost layer of the pericarp is the endocarp. This surrounds a variable number of carpels, shaped as radial segments. The seeds, if present, develop inside the carpels. The space inside each segment is a locule filled with juice vesicles, or pulp. From the endocarp, string-like "hairs" extend into the locules, which provide nourishment to the fruit as it develops. The genus is commercially important with cultivars of many species grown for their fruit. Some cultivars have been developed to be easy to peel and seedless, meaning they are parthenocarpic.
The fragrance of citrus fruits is conferred by flavonoids and limonoids in the rind. The flavonoids include various flavanones and flavones. The carpels are juicy; they contain a high quantity of citric acid, which with other organic acids including ascorbic acid (vitamin C) give them their characteristic sharp taste. Citrus fruits are diverse in size and shape, as well as in color and flavor, reflecting their biochemistry; for instance, grapefruit is made bitter-tasting by a flavanone, naringin.
Cultivation
Most commercial citrus cultivation uses trees produced by grafting the desired fruiting cultivars onto rootstocks selected for disease resistance and hardiness. The trees are not generally frost hardy. They thrive in a consistently sunny, humid environment with fertile soil and adequate water.
The colour of citrus fruits only develops in climates with a (diurnal) cool winter. In tropical regions with no winter at all, citrus fruits remain green until maturity, hence the tropical "green oranges". The terms 'ripe' and 'mature' are widely used synonymously, but they mean different things. A mature fruit is one that has completed its growth phase. Ripening is the sequence of changes within the fruit from maturity to the beginning of decay. These changes involve the conversion of starches to sugars, a decrease in acids, softening, and a change in the fruit's colour. Citrus fruits are non-climacteric and respiration slowly declines and the production and release of ethylene is gradual.
Production
According to the UN Food and Agriculture Organization, world production of all citrus fruits in 2016 was 124 million tonnes, with about half of this production as oranges. At US $15.2 billion equivalent in 2018, citrus trade makes up nearly half of the world fruit trade, which was US$32.1 billion that year. According to the United Nations Conference on Trade and Development, citrus production grew during the early 21st century mainly by the increase in cultivation areas, improvements in transportation and packaging, rising incomes and consumer preference for healthy foods. In 2019–20, world production of oranges was estimated to be 47.5 million tonnes, led by Brazil, Mexico, the European Union, and China as the largest producers.
Pests and diseases
Among the diseases of citrus plantations are citrus black spot (a fungus), citrus canker (a bacterium), citrus greening (a bacterium, spread by an insect pest), and sweet orange scab (a fungus, Elsinöe australis). Citrus plants are liable to infestation by ectoparasites which act as vectors to plant diseases: for example, aphids transmit the damaging citrus tristeza virus, while the aphid-like Asian citrus psyllid can carry the bacterium which causes the serious citrus greening disease. This threatens production in Florida, California, and worldwide. Citrus groves are attacked by parasitic Nematodes including citrus (Tylenchulus semipenetrans) and sheath nematodes (Hemicycliophora spp.).
Deficiency diseases
Citrus plants can develop the deficiency condition chlorosis, characterized by yellowing leaves. The condition is often caused by an excessively high pH (alkaline soil), which prevents the plant from absorbing nutrients such as iron, magnesium, and zinc needed to produce chlorophyll.
Effects on humans
Some Citrus species contain significant amounts of furanocoumarins. In humans, some of these act as strong photosensitizers when applied topically to the skin, while others interact with medications when taken orally in the grapefruit juice effect. Due to the photosensitizing effects of certain furanocoumarins, some Citrus species cause phytophotodermatitis, a potentially severe skin inflammation resulting from contact with a light-sensitizing botanical agent followed by exposure to ultraviolet light. In Citrus species, the primary photosensitizing agent appears to be bergapten, a linear furanocoumarin derived from psoralen. This claim has been confirmed for lime and bergamot. In particular, bergamot essential oil has a higher concentration of bergapten (3–3.6 g/kg) than any other Citrus-based essential oil.
A systematic review indicates that citrus fruit consumption is associated with a 10% reduction of risk for developing breast cancer.
Uses
Culinary
Many citrus fruits, such as oranges, tangerines, grapefruits, and clementines, are generally eaten fresh. They are typically peeled and can be easily split into segments. Grapefruit is more commonly halved and eaten out of the skin with a spoon. Lemonade is a popular beverage prepared by diluting the juice and adding sugar. Lemon juice is mixed in salad dressings and squeezed over fruit salad to stop it from turning brown: its acidity suppresses oxidation by polyphenol oxidase enzymes.
A variety of flavours can be derived from different parts and treatments of citrus fruits. The colourful outer skin of some citrus fruits, known as zest, is used as a flavouring in cooking.
The whole of the bitter orange (and sometimes other citrus fruits) including the peel with its essential oils is cooked with sugar to make marmalade.
As ornamental plants
By the 17th century, orangeries were added to great houses in Europe, both to enable the fruit to be grown locally and for prestige, as seen in the Versailles Orangerie. Some modern hobbyists grow dwarf citrus in containers or greenhouses in areas where the weather is too cold to grow it outdoors; Citrofortunella hybrids have good cold resistance.
In art and culture
Lemons appear in paintings, pop art, and novels. A wall painting in the tomb of Nakht in 15th century BC Egypt depicts a woman in a festival, holding a lemon. In the 17th century, Giovanna Garzoni painted a Still Life with Bowl of Citrons, the fruits still attached to leafy flowering twigs, with a wasp on one of the fruits. The impressionist Edouard Manet depicted a lemon on a pewter plate. In modern art, Arshile Gorky painted Still Life with Lemons in the 1930s.
Citrus fruits "were the clear status symbols of the nobility in the ancient Mediterranean", according to the paleoethnobotanist Dafna Langgut. In Louisa May Alcott's 1868 novel Little Women, the character Amy March states that "It's nothing but limes now, for everyone is sucking them in their desks in schooltime, and trading them off for pencils, bead rings, paper dolls, or something else… If one girl likes another, she gives her a lime; if she’s mad with her, she eats one before her face, and doesn’t offer even a suck."
| Biology and health sciences | Sapindales | null |
52202 | https://en.wikipedia.org/wiki/Magic%20square | Magic square | In mathematics, especially historical and recreational mathematics, a square array of numbers, usually positive integers, is called a magic square if the sums of the numbers in each row, each column, and both main diagonals are the same. The "order" of the magic square is the number of integers along one side (n), and the constant sum is called the "magic constant". If the array includes just the positive integers , the magic square is said to be "normal". Some authors take "magic square" to mean "normal magic square".
Magic squares that include repeated entries do not fall under this definition and are referred to as "trivial". Some well-known examples, including the Sagrada Família magic square and the Parker square, are trivial in this sense. When all the rows and columns but not both diagonals sum to the magic constant, this gives a semimagic square (sometimes called orthomagic square).
The mathematical study of a magic square typically deals with its construction, classification, and enumeration. Although completely general methods for producing all the magic squares of all orders do not exist, historically three general techniques have been discovered: by bordering method, by making composite magic squares, and by adding two preliminary squares. There are also more specific strategies like the continuous enumeration method that reproduces specific patterns. Magic squares are generally classified according to their order n as: odd if n is odd, evenly even (also referred to as "doubly even") if n is a multiple of 4, oddly even (also known as "singly even") if n is any other even number. This classification is based on different techniques required to construct odd, evenly even, and oddly even squares. Beside this, depending on further properties, magic squares are also classified as associative magic squares, pandiagonal magic squares, most-perfect magic squares, and so on. More challengingly, attempts have also been made to classify all the magic squares of a given order as transformations of a smaller set of squares. Except for n ≤ 5, the enumeration of higher order magic squares is still an open challenge. The enumeration of most-perfect magic squares of any order was only accomplished in the late 20th century.
Magic squares have a long history, dating back to at least 190 BCE in China. At various times they have acquired occult or mythical significance, and have appeared as symbols in works of art. In modern times they have been generalized a number of ways, including using extra or different constraints, multiplying instead of adding cells, using alternate shapes or more than two dimensions, and replacing numbers with shapes and addition with geometric operations.
History
The third-order magic square was known to Chinese mathematicians as early as 190 BCE, and explicitly given by the first century of the common era. The first dateable instance of the fourth-order magic square occurred in 587 CE in India. Specimens of magic squares of order 3 to 9 appear in an encyclopedia from Baghdad , the Encyclopedia of the Brethren of Purity (Rasa'il Ikhwan al-Safa). By the end of 12th century, the general methods for constructing magic squares were well established. Around this time, some of these squares were increasingly used in conjunction with magic letters, as in Shams Al-ma'arif, for occult purposes. In India, all the fourth-order pandiagonal magic squares were enumerated by Narayana in 1356. Magic squares were made known to Europe through translation of Arabic sources as occult objects during the Renaissance, and the general theory had to be re-discovered independent of prior developments in China, India, and Middle East. Also notable are the ancient cultures with a tradition of mathematics and numerology that did not discover the magic squares: Greeks, Babylonians, Egyptians, and Pre-Columbian Americans.
Chinese
While ancient references to the pattern of even and odd numbers in the 3×3 magic square appear in the I Ching, the first unequivocal instance of this magic square appears in the chapter called Mingtang (Bright Hall) of a 1st-century book Da Dai Liji (Record of Rites by the Elder Dai), which purported to describe ancient Chinese rites of the Zhou dynasty.
These numbers also occur in a possibly earlier mathematical text called Shushu jiyi (Memoir on Some Traditions of Mathematical Art), said to be written in 190 BCE. This is the earliest appearance of a magic square on record; and it was mainly used for divination and astrology. The 3×3 magic square was referred to as the "Nine Halls" by earlier Chinese mathematicians. The identification of the 3×3 magic square to the legendary Luoshu chart was only made in the 12th century, after which it was referred to as the Luoshu square. The oldest surviving Chinese treatise that displays magic squares of order larger than 3 is Yang Hui's Xugu zheqi suanfa (Continuation of Ancient Mathematical Methods for Elucidating the Strange) written in 1275. The contents of Yang Hui's treatise were collected from older works, both native and foreign; and he only explains the construction of third and fourth-order magic squares, while merely passing on the finished diagrams of larger squares. He gives a magic square of order 3, two squares for each order of 4 to 8, one of order nine, and one semi-magic square of order 10. He also gives six magic circles of varying complexity.
The above magic squares of orders 3 to 9 are taken from Yang Hui's treatise, in which the Luo Shu principle is clearly evident. The order 5 square is a bordered magic square, with central 3×3 square formed according to Luo Shu principle. The order 9 square is a composite magic square, in which the nine 3×3 sub squares are also magic. After Yang Hui, magic squares frequently occur in Chinese mathematics such as in Ding Yidong's Dayan suoyin (), Cheng Dawei's Suanfa tongzong (1593), Fang Zhongtong's Shuduyan (1661) which contains magic circles, cubes and spheres, Zhang Chao's Xinzhai zazu (), who published China's first magic square of order ten, and lastly Bao Qishou's Binaishanfang ji (), who gave various three dimensional magic configurations. However, despite being the first to discover the magic squares and getting a head start by several centuries, the Chinese development of the magic squares are much inferior compared to the Indian, Middle Eastern, or European developments. The high point of Chinese mathematics that deals with the magic squares seems to be contained in the work of Yang Hui; but even as a collection of older methods, this work is much more primitive, lacking general methods for constructing magic squares of any order, compared to a similar collection written around the same time by the Byzantine scholar Manuel Moschopoulos. This is possibly because of the Chinese scholars' enthralment with the Lo Shu principle, which they tried to adapt to solve higher squares; and after Yang Hui and the fall of Yuan dynasty, their systematic purging of the foreign influences in Chinese mathematics.
Japan
Japan and China have similar mathematical traditions and have repeatedly influenced each other in the history of magic squares. The Japanese interest in magic squares began after the dissemination of Chinese works—Yang Hui's Suanfa and Cheng Dawei's Suanfa tongzong—in the 17th century, and as a result, almost all the wasans devoted their time to its study.
In the 1660 edition of Ketsugi-sho, Isomura Kittoku gave both odd and even ordered bordered magic squares as well as magic circles; while the 1684 edition of the same book contained a large section on magic squares, demonstrating that he had a general method for constructing bordered magic squares. In Jinko-ki (1665) by Muramatsu Kudayu Mosei, both magic squares and magic circles are displayed. The largest square Mosei constructs is of 19th order. Various magic squares and magic circles were also published by Nozawa Teicho in Dokai-sho (1666), Sato Seiko in Kongenki (1666), and Hosino Sanenobu in Ko-ko-gen Sho (1673). One of Seki Takakazu's Seven Books (Hojin Yensan) (1683) is devoted completely to magic squares and circles. This is the first Japanese book to give a general treatment of magic squares in which the algorithms for constructing odd, singly even and doubly even bordered magic squares are clearly described. In 1694 and 1695, Yueki Ando gave different methods to create the magic squares and displayed squares of order 3 to 30. A fourth-order magic cube was constructed by Yoshizane Tanaka (1651–1719) in Rakusho-kikan (1683). The study of magic squares was continued by Seki's pupils, notably by Katahiro Takebe, whose squares were displayed in the fourth volume of Ichigen Kappo by Shukei Irie, Yoshisuke Matsunaga in Hojin-Shin-jutsu, Yoshihiro Kurushima in Kyushi Iko who rediscovered a method to produce the odd squares given by Agrippa, and Naonobu Ajima. Thus by the beginning of the 18th century, the Japanese mathematicians were in possession of methods to construct magic squares of arbitrary order. After this, attempts at enumerating the magic squares was initiated by Nushizumi Yamaji.
India
The 3×3 magic square first appears in India in Gargasamhita by Garga, who recommends its use to pacify the nine planets (navagraha). The oldest version of this text dates from 100 CE, but the passage on planets could not have been written earlier than 400 CE. The first dateable instance of 3×3 magic square in India occur in a medical text Siddhayog () by Vrnda, which was prescribed to women in labor in order to have easy delivery.
The oldest dateable fourth order magic square in the world is found in an encyclopaedic work written by Varahamihira around 587 CE called Brhat Samhita. The magic square is constructed for the purpose of making perfumes using 4 substances selected from 16 different substances. Each cell of the square represents a particular ingredient, while the number in the cell represents the proportion of the associated ingredient, such that the mixture of any four combination of ingredients along the columns, rows, diagonals, and so on, gives the total volume of the mixture to be 18. Although the book is mostly about divination, the magic square is given as a matter of combinatorial design, and no magical properties are attributed to it. The special features of this magic square were commented on by Bhattotpala ()
The square of Varahamihira as given above has sum of 18. Here the numbers 1 to 8 appear twice in the square. It is a pan-diagonal magic square. Four different magic squares can be obtained by adding 8 to one of the two sets of 1 to 8 sequence. The sequence is selected such that the number 8 is added exactly twice in each row, each column and each of the main diagonals. One of the possible magic squares shown in the right side. This magic square is remarkable in that it is a 90 degree rotation of a magic square that appears in the 13th century Islamic world as one of the most popular magic squares.
The construction of 4th-order magic square is detailed in a work titled Kaksaputa, composed by the alchemist Nagarjuna around 10th century CE. All of the squares given by Nagarjuna are 4×4 magic squares, and one of them is called Nagarjuniya after him. Nagarjuna gave a method of constructing 4×4 magic square using a primary skeleton square, given an odd or even magic sum. The Nagarjuniya square is given below, and has the sum total of 100.
The Nagarjuniya square is a pan-diagonal magic square. The Nagarjuniya square is made up of two arithmetic progressions starting from 6 and 16 with eight terms each, with a common difference between successive terms as 4. When these two progressions are reduced to the normal progression of 1 to 8, the adjacent square is obtained.
Around 12th-century, a 4×4 magic square was inscribed on the wall of Parshvanath temple in Khajuraho, India. Several Jain hymns teach how to make magic squares, although they are undateable.
As far as is known, the first systematic study of magic squares in India was conducted by Thakkar Pheru, a Jain scholar, in his Ganitasara Kaumudi (c. 1315). This work contains a small section on magic squares which consists of nine verses. Here he gives a square of order four, and alludes to its rearrangement; classifies magic squares into three (odd, evenly even, and oddly even) according to its order; gives a square of order six; and prescribes one method each for constructing even and odd squares. For the even squares, Pheru divides the square into component squares of order four, and puts the numbers into cells according to the pattern of a standard square of order four. For odd squares, Pheru gives the method using horse move or knight's move. Although algorithmically different, it gives the same square as the De la Loubere's method.
The next comprehensive work on magic squares was taken up by Narayana Pandit, who in the fourteenth chapter of his Ganita Kaumudi (1356) gives general methods for their construction, along with the principles governing such constructions. It consists of 55 verses for rules and 17 verses for examples. Narayana gives a method to construct all the pan-magic squares of fourth order using knight's move; enumerates the number of pan-diagonal magic squares of order four, 384, including every variation made by rotation and reflection; three general methods for squares having any order and constant sum when a standard square of the same order is known; two methods each for constructing evenly even, oddly even, and of squares when the sum is given. While Narayana describes one older method for each species of square, he claims the method of superposition for evenly even and odd squares and a method of interchange for oddly even squares to be his own invention. The superposition method was later re-discovered by De la Hire in Europe. In the last section, he conceives of other figures, such as circles, rectangles, and hexagons, in which the numbers may be arranged to possess properties similar to those of magic squares. Below are some of the magic squares constructed by Narayana:
The order 8 square is interesting in itself since it is an instance of the most-perfect magic square. Incidentally, Narayana states that the purpose of studying magic squares is to construct yantra, to destroy the ego of bad mathematicians, and for the pleasure of good mathematicians. The subject of magic squares is referred to as bhadraganita and Narayana states that it was first taught to men by god Shiva.
Middle East, North Africa, Muslim Iberia
Although the early history of magic squares in Persia and Arabia is not known, it has been suggested that they were known in pre-Islamic times. It is clear, however, that the study of magic squares was common in medieval Islam, and it was thought to have begun after the introduction of chess into the region. The first dateable appearance of a magic square of order 3 occurs in Jābir ibn Hayyān's (fl. c. 721 – c. 815) Kitab al-mawazin al-Saghir (The Small Book of Balances) where the magic square and its related numerology is associated with alchemy. While it is known that treatises on magic squares were written in the 9th century, the earliest extant treaties date from the 10th-century: one by Abu'l-Wafa al-Buzjani () and another by Ali b. Ahmad al-Antaki (). These early treatises were purely mathematical, and the Arabic designation for magic squares used is wafq al-a'dad, which translates as harmonious disposition of the numbers. By the end of 10th century, the two treatises by Buzjani and Antaki makes it clear that the Middle Eastern mathematicians had understood how to construct bordered squares of any order as well as simple magic squares of small orders (n ≤ 6) which were used to make composite magic squares. A specimen of magic squares of orders 3 to 9 devised by Middle Eastern mathematicians appear in an encyclopedia from Baghdad , the Rasa'il Ikhwan al-Safa (the Encyclopedia of the Brethren of Purity). The squares of order 3 to 7 from Rasa'il are given below:
The 11th century saw the finding of several ways to construct simple magic squares for odd and evenly-even orders; the more difficult case of evenly-odd case (n = 4k + 2) was solved by Ibn al-Haytham with k even (c. 1040), and completely by the beginning of 12th century, if not already in the latter half of the 11th century. Around the same time, pandiagonal squares were being constructed. Treaties on magic squares were numerous in the 11th and 12th century. These later developments tended to be improvements on or simplifications of existing methods. From the 13th century, magic squares were increasingly put to occult purposes. However, much of these later texts written for occult purposes merely depict certain magic squares and mention their attributes, without describing their principle of construction, with only some authors keeping the general theory alive. One such occultist was the Algerian Ahmad al-Buni (c. 1225), who gave general methods on constructing bordered magic squares; some others were the 17th century Egyptian Shabramallisi and the 18th century Nigerian al-Kishnawi.
The magic square of order three was described as a child-bearing charm since its first literary appearances in the alchemical works of Jābir ibn Hayyān (fl. c. 721 – c. 815) and al-Ghazālī (1058–1111) and it was preserved in the tradition of the planetary tables. The earliest occurrence of the association of seven magic squares to the virtues of the seven heavenly bodies appear in Andalusian scholar Ibn Zarkali's (known as Azarquiel in Europe) (1029–1087) Kitāb tadbīrāt al-kawākib (Book on the Influences of the Planets). A century later, the Algerian scholar Ahmad al-Buni attributed mystical properties to magic squares in his highly influential book Shams al-Ma'arif (The Book of the Sun of Gnosis and the Subtleties of Elevated Things), which also describes their construction. This tradition about a series of magic squares from order three to nine, which are associated with the seven planets, survives in Greek, Arabic, and Latin versions. There are also references to the use of magic squares in astrological calculations, a practice that seems to have originated with the Arabs.
Latin Europe
Unlike in Persia and Arabia, better documentation exists of how the magic squares were transmitted to Europe. Around 1315, influenced by Arab sources, the Greek Byzantine scholar Manuel Moschopoulos wrote a mathematical treatise on the subject of magic squares, leaving out the mysticism of his Middle Eastern predecessors, where he gave two methods for odd squares and two methods for evenly even squares. Moschopoulos was essentially unknown to the Latin Europe until the late 17th century, when Philippe de la Hire rediscovered his treatise in the Royal Library of Paris. However, he was not the first European to have written on magic squares; and the magic squares were disseminated to rest of Europe through Spain and Italy as occult objects. The early occult treaties that displayed the squares did not describe how they were constructed. Thus the entire theory had to be rediscovered.
Magic squares had first appeared in Europe in Kitāb tadbīrāt al-kawākib (Book on the Influences of the Planets) written by Ibn Zarkali of Toledo, Al-Andalus, as planetary squares by 11th century. The magic square of three was discussed in numerological manner in early 12th century by Jewish scholar Abraham ibn Ezra of Toledo, which influenced later Kabbalists. Ibn Zarkali's work was translated as Libro de Astromagia in the 1280s, due to Alfonso X of Castille. In the Alfonsine text, magic squares of different orders are assigned to the respective planets, as in the Islamic literature; unfortunately, of all the squares discussed, the Mars magic square of order five is the only square exhibited in the manuscript.
Magic squares surface again in Florence, Italy in the 14th century. A 6×6 and a 9×9 square are exhibited in a manuscript of the Trattato d'Abbaco (Treatise of the Abacus) by Paolo Dagomari. It is interesting to observe that Paolo Dagomari, like Pacioli after him, refers to the squares as a useful basis for inventing mathematical questions and games, and does not mention any magical use. Incidentally, though, he also refers to them as being respectively the Sun's and the Moon's squares, and mentions that they enter astrological calculations that are not better specified. As said, the same point of view seems to motivate the fellow Florentine Luca Pacioli, who describes 3×3 to 9×9 squares in his work De Viribus Quantitatis by the end of 15th century.
Europe after 15th century
The planetary squares had disseminated into northern Europe by the end of 15th century. For instance, the Cracow manuscript of Picatrix from Poland displays magic squares of orders 3 to 9. The same set of squares as in the Cracow manuscript later appears in the writings of Paracelsus in Archidoxa Magica (1567), although in highly garbled form. In 1514 Albrecht Dürer immortalized a 4×4 square in his famous engraving Melencolia I. Paracelsus' contemporary Heinrich Cornelius Agrippa von Nettesheim published his famous three volume book De occulta philosophia in 1531, where he devoted Chapter 22 of Book II to the planetary squares shown below. The same set of squares given by Agrippa reappear in 1539 in Practica Arithmetice by Girolamo Cardano, where he explains the construction of the odd ordered squares using "diamond method", which was later reproduced by Bachet. The tradition of planetary squares was continued into the 17th century by Athanasius Kircher in Oedipi Aegyptici (1653). In Germany, mathematical treaties concerning magic squares were written in 1544 by Michael Stifel in Arithmetica Integra, who rediscovered the bordered squares, and Adam Riese, who rediscovered the continuous numbering method to construct odd ordered squares published by Agrippa. However, due to the religious upheavals of that time, these work were unknown to the rest of Europe.
In 1624 France, Claude Gaspard Bachet described the "diamond method" for constructing Agrippa's odd ordered squares in his book Problèmes Plaisants. During 1640 Bernard Frenicle de Bessy and Pierre Fermat exchanged letters on magic squares and cubes, and in one of the letters Fermat boasts of being able to construct 1,004,144,995,344 magic squares of order 8 by his method. An early account on the construction of bordered squares was given by Antoine Arnauld in his Nouveaux éléments de géométrie (1667). In the two treatise Des quarrez ou tables magiques and Table générale des quarrez magiques de quatre de côté, published posthumously in 1693, twenty years after his death, Bernard Frenicle de Bessy demonstrated that there were exactly 880 distinct magic squares of order four. Frenicle gave methods to construct magic square of any odd and even order, where the even ordered squares were constructed using borders. He also showed that interchanging rows and columns of a magic square produced new magic squares. In 1691, Simon de la Loubère described the Indian continuous method of constructing odd ordered magic squares in his book Du Royaume de Siam, which he had learned while returning from a diplomatic mission to Siam, which was faster than Bachet's method. In an attempt to explain its working, de la Loubere used the primary numbers and root numbers, and rediscovered the method of adding two preliminary squares. This method was further investigated by Abbe Poignard in Traité des quarrés sublimes (1704), by Philippe de La Hire in Mémoires de l'Académie des Sciences for the Royal Academy (1705), and by Joseph Sauveur in Construction des quarrés magiques (1710). Concentric bordered squares were also studied by De la Hire in 1705, while Sauveur introduced magic cubes and lettered squares, which was taken up later by Euler in 1776, who is often credited for devising them. In 1750 d'Ons-le-Bray rediscovered the method of constructing doubly even and singly even squares using bordering technique; while in 1767 Benjamin Franklin published a semi-magic square that had the properties of eponymous Franklin square. By this time the earlier mysticism attached to the magic squares had completely vanished, and the subject was treated as a part of recreational mathematics.
In the 19th century, Bernard Violle gave a comprehensive treatment of magic squares in his three volume Traité complet des carrés magiques (1837–1838), which also described magic cubes, parallelograms, parallelopipeds, and circles. Pandiagonal squares were extensively studied by Andrew Hollingworth Frost, who learned it while in the town of Nasik, India, (thus calling them Nasik squares) in a series of articles: On the knight's path (1877), On the General Properties of Nasik Squares (1878), On the General Properties of Nasik Cubes (1878), On the construction of Nasik Squares of any order (1896). He showed that it is impossible to have normal singly-even pandiagonal magic squares. Frederick A.P. Barnard constructed inlaid magic squares and other three dimensional magic figures like magic spheres and magic cylinders in Theory of magic squares and of magic cubes (1888). In 1897, Emroy McClintock published On the most perfect form of magic squares, coining the words pandiagonal square and most perfect square, which had previously been referred to as perfect, or diabolic, or Nasik.
Some famous magic squares
Luo Shu magic square
Legends dating from as early as 650 BCE tell the story of the Lo Shu (洛書) or "scroll of the river Lo". According to the legend, there was at one time in ancient China a huge flood. While the great king Yu was trying to channel the water out to sea, a turtle emerged from it with a curious pattern on its shell: a 3×3 grid in which circular dots of numbers were arranged, such that the sum of the numbers in each row, column and diagonal was the same: 15. According to the legend, thereafter people were able to use this pattern in a certain way to control the river and protect themselves from floods. The Lo Shu Square, as the magic square on the turtle shell is called, is the unique normal magic square of order three in which 1 is at the bottom and 2 is in the upper right corner. Every normal magic square of order three is obtained from the Lo Shu by rotation or reflection.
Magic square in Parshavnath temple
There is a well-known 12th-century 4×4 normal magic square inscribed on the wall of the Parshvanath temple in Khajuraho, India.
This is known as the Chautisa Yantra (Chautisa, 34; Yantra, lit. "device"), since its magic sum is 34. It is one of the three 4×4 pandiagonal magic squares and is also an instance of the most-perfect magic square. The study of this square led to the appreciation of pandiagonal squares by European mathematicians in the late 19th century. Pandiagonal squares were referred to as Nasik squares or Jain squares in older English literature.
Albrecht Dürer's magic square
The order four normal magic square Albrecht Dürer immortalized in his 1514 engraving Melencolia I, referred to above, is believed to be the first seen in European art. The square associated with Jupiter appears as a talisman used to drive away melancholy. It is very similar to Yang Hui's square, which was created in China about 250 years before Dürer's time. As with every order 4 normal magic square, the magic sum is 34. But in the Durer square this sum is also found
in each of the quadrants, in the center four squares, and in the corner squares (of the 4×4 as well as the four contained 3×3 grids). This sum can also be found in the four outer numbers clockwise from the corners (3+8+14+9) and likewise the four counter-clockwise (the locations of four queens in the two solutions of the 4 queens puzzle), the two sets of four symmetrical numbers (2+8+9+15 and 3+5+12+14), the sum of the middle two entries of the two outer columns and rows (5+9+8+12 and 3+2+15+14), and in four kite or cross shaped quartets (3+5+11+15, 2+10+8+14, 3+9+7+15, and 2+6+12+14). The two numbers in the middle of the bottom row give the date of the engraving: 1514. It has been speculated that the numbers 4,1 bordering the publication date correspond to Durer's initials D,A. But if that had been his intention, he could have inverted the order of columns 1 and 4 to achieve "A1514D" without compromising the square's properties.
Dürer's magic square can also be extended to a magic cube.
Sagrada Família magic square
The Passion façade of the Sagrada Família church in Barcelona, conceptualized by Antoni Gaudí and designed by sculptor Josep Subirachs, features a trivial order 4 magic square: The magic constant of the square is 33, the age of Jesus at the time of the Passion. Structurally, it is very similar to the Melancholia magic square, but it has had the numbers in four of the cells reduced by 1.
Trivial squares such as this one are not generally mathematically interesting and only have historical significance. Lee Sallows has pointed out that, due to Subirachs's ignorance of magic square theory, the renowned sculptor made a needless blunder, and supports this assertion by giving several examples of non-trivial 4×4 magic squares showing the desired magic constant of 33.
Similarly to Dürer's magic square, the Sagrada Familia's magic square can also be extended to a magic cube.
Parker square
The Parker square, named after recreational mathematician Matt Parker, is an attempt to create a 33 magic square of squares — a prized unsolved problem since Euler. The Parker square is a trivial semimagic square since it uses some numbers more than once, and the diagonal sums to , not as for all the other rows and columns, and the other diagonal. The Parker square became popular in mathematical culture. The Parker square became a "mascot for people who give it a go, but ultimately fall short".
Gardner square
The Gardner square, named after recreational mathematician Martin Gardner, similar to the Parker square,
is given as a problem to determine a, b, c and d.
This solution for a = 74, b = 113, c = 94 and d = 97 gives a semimagic square; the diagonal sums to , not as for all the other rows and columns, and the other diagonal.
Properties of magic squares
Magic constant
The constant that is the sum of any row, or column, or diagonal is called the magic constant or magic sum, M. Every normal magic square has a constant dependent on the order , calculated by the formula . This can be demonstrated by noting that the sum of is . Since the sum of each row is , the sum of rows is , which when divided by the order yields the magic constant as . For normal magic squares of orders n = 3, 4, 5, 6, 7, and 8, the magic constants are, respectively: 15, 34, 65, 111, 175, and 260 (sequence A006003 in the OEIS).
Magic square of order 1 is trivial
The 1×1 magic square, with only one cell containing the number 1, is called trivial, because it is typically not under consideration when discussing magic squares; but it is indeed a magic square by definition, if a single cell is regarded as a square of order one.
Magic square of order 2 cannot be constructed
Normal magic squares of all sizes can be constructed except 2×2 (that is, where order n = 2).
Center of mass
If the numbers in the magic square are seen as masses located in various cells, then the center of mass of a magic square coincides with its geometric center.
Moment of inertia
The moment of inertia of a magic square has been defined as the sum over all cells of the number in the cell times the squared distance from the center of the cell to the center of the square; here the unit of measurement is the width of one cell. (Thus for example a corner cell of a 3×3 square has a distance of a non-corner edge cell has a distance of 1, and the center cell has a distance of 0.) Then all magic squares of a given order have the same moment of inertia as each other. For the order-3 case the moment of inertia is always 60, while for the order-4 case the moment of inertia is always 340. In general, for the n×n case the moment of inertia is
Birkhoff–von Neumann decomposition
Dividing each number of the magic square by the magic constant will yield a doubly stochastic matrix, whose row sums and column sums equal to unity. However, unlike the doubly stochastic matrix, the diagonal sums of such matrices will also equal to unity. Thus, such matrices constitute a subset of doubly stochastic matrix. The Birkhoff–von Neumann theorem states that for any doubly stochastic matrix , there exists real numbers , where and permutation matrices such that
This representation may not be unique in general. By Marcus-Ree theorem, however, there need not be more than terms in any decomposition. Clearly, this decomposition carries over to magic squares as well, since a magic square can be recovered from a doubly stochastic matrix by multiplying it by the magic constant.
Classification of magic squares
While the classification of magic squares can be done in many ways, some useful categories are given below. An n×n square array of integers 1, 2, ..., n2 is called:
Semi-magic square when its rows and columns sum to give the magic constant.
Simple magic square when its rows, columns, and two diagonals sum to give magic constant and no more. They are also known as ordinary magic squares or normal magic squares.
Self-complementary magic square when it is a magic square which when complemented (i.e. each number subtracted from n2 + 1) will give a rotated or reflected version of the original magic square.
Associative magic square when it is a magic square with a further property that every number added to the number equidistant, in a straight line, from the center gives n2 + 1. They are also called symmetric magic squares. Associative magic squares do not exist for squares of singly even order. All associative magic square are self-complementary magic squares as well.
Pandiagonal magic square when it is a magic square with a further property that the broken diagonals sum to the magic constant. They are also called panmagic squares, perfect squares, diabolic squares, Jain squares, or Nasik squares. Panmagic squares do not exist for singly even orders. However, singly even non-normal squares can be panmagic.
Ultra magic square when it is both associative and pandiagonal magic square. Ultra magic square exist only for orders n ≥ 5.
Bordered magic square when it is a magic square and it remains magic when the rows and columns on the outer edge are removed. They are also called concentric bordered magic squares if removing a border of a square successively gives another smaller bordered magic square. Bordered magic square do not exist for order 4.
Composite magic square when it is a magic square that is created by "multiplying" (in some sense) smaller magic squares, such that the order of the composite magic square is a multiple of the order of the smaller squares. Such squares can usually be partitioned into smaller non-overlapping magic sub-squares.
Inlaid magic square when it is a magic square inside which a magic sub-square is embedded, regardless of construction technique. The embedded magic sub-squares are themselves referred to as inlays.
Most-perfect magic square when it is a pandiagonal magic square with two further properties (i) each 2×2 subsquare add to 1/k of the magic constant where n = 4k, and (ii) all pairs of integers distant n/2 along any diagonal (major or broken) are complementary (i.e. they sum to n2 + 1). The first property is referred to as compactness, while the second property is referred to as completeness. Most-perfect magic squares exist only for squares of doubly even order. All the pandiagonal squares of order 4 are also most perfect.
Franklin magic square when it is a doubly even magic square with three further properties (i) every bent diagonal adds to the magic constant, (ii) every half row and half column starting at an outside edge adds to half the magic constant, and (iii) the square is compact.
Multimagic square when it is a magic square that remains magic even if all its numbers are replaced by their k-th power for 1 ≤ k ≤ P. They are also known as P-multimagic square or satanic squares. They are also referred to as bimagic squares, trimagic squares, tetramagic squares, and pentamagic squares when the value of P is 2, 3, 4, and 5 respectively.
Enumeration of magic squares
Low-order squares
There is only one (trivial) magic square of order 1 and no magic square of order 2. As mentioned above, the set of normal squares of order three constitutes a single equivalence class-all equivalent to the Lo Shu square. Thus there is basically just one normal magic square of order 3.
The number of different n × n magic squares for n from 1 to 6, not counting rotations and reflections is:
1, 0, 1, 880, 275305224, 17753889197660635632.
The number for n = 6 had previously been estimated to be
Magic tori
Cross-referenced to the above sequence, a new classification enumerates the magic tori that display these magic squares. The number of magic tori of order n from 1 to 5, is:
1, 0, 1, 255, 251449712 .
Higher-order squares and tori
The number of distinct normal magic squares rapidly increases for higher orders.
The 880 magic squares of order 4 are displayed on 255 magic tori of order 4 and the 275,305,224 squares of order 5 are displayed on 251,449,712 magic tori of order 5. The numbers of magic tori and distinct normal squares are not yet known for orders beyond 5 and 6, respectively.
Algorithms tend to only generate magic squares of a certain type or classification, making counting all possible magic squares quite difficult. Since traditional counting methods have proven unsuccessful, statistical analysis using the Monte Carlo method has been applied. The basic principle applied to magic squares is to randomly generate n × n matrices of elements 1 to n2 and check if the result is a magic square. The probability that a randomly generated matrix of numbers is a magic square is then used to approximate the number of magic squares.
More intricate versions of the Monte Carlo method, such as the exchange Monte Carlo, and Monte Carlo backtracking have produced even more accurate estimations. Using these methods it has been shown that the probability of magic squares decreases rapidly as n increases. Using fitting functions give the curves seen to the right.
Transformations that preserve the magic property
For any magic square
The sum of any two magic squares of the same order by matrix addition is a magic square.
A magic square remains magic when all of its numbers undergo the same linear transformation (i.e., a function of the form ). For example, a magic square remains magic when its numbers are multiplied by any constant. Moreover, a magic square remains magic when a constant is added or subtracted to its numbers, or if its numbers are subtracted from a constant. In particular, if every element in a normal magic square of order is subtracted from , the complement of the original square is obtained. In the example below, each element of the magic square on the left is subtracted from 17 to obtain the complement magic square on the right.
A magic square remains magic when transformed by any element of , the symmetry group of a square (see ). Every combination of one or more rotations of 90 degrees, reflections, or both produce eight trivially distinct squares which are generally considered equivalent. The eight such squares are said to make up a single equivalence class. The eight equivalent magic squares for the 3×3 magic square are shown below:
A magic square of order remains magic when both its rows and columns are symmetrically permuted by such that for . Every permutation of the rows or columns preserves all row and column sums, but generally not the two diagonal sums. If the same permutation is applied to both the rows and columns, then diagonal element in row and column is mapped to row and column which is on the same diagonal; therefore, applying the same permutation to rows and columns preserves the main (upper left to lower right) diagonal sum. If the permutation is symmetric as described, then the diagonal element in row and column is mapped to row and column which is on the same diagonal; therefore, applying the same symmetric permutation to both rows and columns preserves both diagonal sums. For even , there are such symmetric permutations, and for odd. In the example below, the original magic square on the left has its rows and columns symmetrically permuted by resulting in the magic square on the right.
A magic square of order remains magic when rows and are exchanged and columns and are exchanged because this is a symmetric permutation of the form described above. In the example below, the square on the right is obtained by interchanging the 1st and 4th rows and columns of the original square on the left.
A magic square of order remains magic when rows and are exchanged, rows and are exchanged, columns and are exchanged, and columns and are exchanged where because this is another symmetric permutation of the form described above. In the example below, the left square is the original square, while the right square is the new square obtained by this transformation. In the middle square, rows 1 and 2 and rows 3 and 4 have been swapped. The final square on the right is obtained by interchanging columns 1 and 2 and columns 3 and 4 of the middle square. In this particular example, this transform rotates the quadrants 180 degrees. The middle square is also magic because the original square is associative.
A magic square remains magic when its quadrants are diagonally interchanged because this is another symmetric permutation of the form described above. For even-order , permute the rows and columns by permutation where for , and for . For odd-order , permute rows and columns by permutation where for , and for . For odd ordered square, the halves of the central row and column are also interchanged. Examples for order 4 and 5 magic squares are given below:
For associative magic squares
An associative magic square remains associative when two rows or columns equidistant from the center are interchanged. For an even square, there are n/2 pairs of rows or columns that can be interchanged; thus 2n/2 × 2n/2 = 2n equivalent magic squares by combining such interchanges can be obtained. For odd square, there are (n - 1)/2 pairs of rows or columns that can be interchanged; and 2n-1 equivalent magic squares obtained by combining such interchanges. Interchanging all the rows flips the square vertically (i.e. reflected along the horizontal axis), while interchanging all the columns flips the square horizontally (i.e. reflected along the vertical axis). In the example below, a 4×4 associative magic square on the left is transformed into a square on the right by interchanging the second and third row, yielding the famous Durer's magic square.
An associative magic square remains associative when two same sided rows (or columns) are interchanged along with corresponding other sided rows (or columns). For an even square, since there are n/2 same sided rows (or columns), there are n(n - 2)/8 pairs of such rows (or columns) that can be interchanged. Thus, 2n(n-2)/8 × 2n(n-2)/8 = 2n(n-2)/4 equivalent magic squares can be obtained by combining such interchanges. For odd square, since there are (n - 1)/2 same sided rows or columns, there are (n - 1)(n - 3)/8 pairs of such rows or columns that can be interchanged. Thus, there are 2(n - 1)(n - 3)/8 × 2(n - 1)(n - 3)/8 = 2(n - 1)(n - 3)/4 equivalent magic squares obtained by combining such interchanges. Interchanging all the same sided rows flips each quadrants of the square vertically, while interchanging all the same sided columns flips each quadrant of the square horizontally. In the example below, the original square is on the left, whose rows 1 and 2 are interchanged with each other, along with rows 3 and 4, to obtain the transformed square on the right.
An associative magic square remains associative when its entries are replaced with corresponding numbers from a set of s arithmetic progressions with the same common difference among r terms, such that r × s = n2, and whose initial terms are also in arithmetic progression, to obtain a non-normal magic square. Here either s or r should be a multiple of n. Let us have s arithmetic progressions given by
where a is the initial term, c is the common difference of the arithmetic progressions, and d is the common difference among the initial terms of each progression. The new magic constant will be
If s = r = n, then follows the simplification
With a = c = 1 and d = n, the usual M = n(n2+1)/2 is obtained. For given M the required a, c, and d can be found by solving the linear Diophantine equation. In the examples below, there are order 4 normal magic squares on the left most side. The second square is a corresponding non-normal magic square with r = 8, s = 2, a = 1, c = 1, and d = 10 such that the new magic constant is M = 38. The third square is an order 5 normal magic square, which is a 90 degree clockwise rotated version of the square generated by De la Loubere method. On the right most side is a corresponding non-normal magic square with a = 4, c = 1, and d = 6 such that the new magic constant is M = 90.
For pan-diagonal magic squares
A pan-diagonal magic square remains a pan-diagonal magic square under cyclic shifting of rows or of columns or both. This allows us to position a given number in any one of the n2 cells of an n order square. Thus, for a given pan-magic square, there are n2 equivalent pan-magic squares. In the example below, the original square on the left is transformed by shifting the first row to the bottom to obtain a new pan-magic square in the middle. Next, the 1st and 2nd column of the middle pan-magic square is circularly shifted to the right to obtain a new pan-magic square on the right.
For bordered magic squares
A bordered magic square remains a bordered magic square after permuting the border cells in the rows or columns, together with their corresponding complementary terms, keeping the corner cells fixed. Since the cells in each row and column of every concentric border can be permuted independently, when the order n ≥ 5 is odd, there are ((n-2)! × (n-4)! × ··· × 3!)2 equivalent bordered squares. When n ≥ 6 is even, there are ((n-2)! × (n-4)! × ··· × 4!)2 equivalent bordered squares. In the example below, a square of order 5 is given whose border row has been permuted and (3!)2 = 36 such equivalent squares can be obtained.
A bordered magic square remains a bordered magic square after each of its concentric borders are independently rotated or reflected with respect to the central core magic square. If there are b borders, then this transform will yield 8b equivalent squares. In the example below of the 5×5 magic square, the border has been rotated 90 degrees anti-clockwise.
For composite magic squares
A composite magic square remains a composite magic square when the embedded magic squares undergo transformations that do not disturb the magic property (e.g. rotation, reflection, shifting of rows and columns, and so on).
Special methods of construction
Over the millennia, many ways to construct magic squares have been discovered. These methods can be classified as general methods and special methods, in the sense that general methods allow us to construct more than a single magic square of a given order, whereas special methods allow us to construct just one magic square of a given order. Special methods are specific algorithms whereas general methods may require some trial-and-error.
Special methods are the most simple ways to construct magic squares. They follow certain algorithms which generate regular patterns of numbers in a square. The correctness of these special methods can be proved using one of the general methods given in later sections. After a magic square has been constructed using a special method, the transformations described in the previous section can be applied to yield further magic squares. Special methods are usually referred to using the name of the author(s) (if known) who described the method, for e.g. De la Loubere's method, Starchey's method, Bachet's method, etc.
Magic squares are believed to exist for all orders, except for order 2. Magic squares can be classified according to their order as odd, doubly even (n divisible by four), and singly even (n even, but not divisible by four). This classification is based on the fact that entirely different techniques need to be employed to construct these different species of squares. Odd and doubly even magic squares are easy to generate; the construction of singly even magic squares is more difficult but several methods exist, including John Horton Conway's LUX method for magic squares and the Strachey method for magic squares.
A method for constructing a magic square of order 3
In the 19th century, Édouard Lucas devised the general formula for order 3 magic squares. Consider the following table made up of positive integers a, b and c:
These nine numbers will be distinct positive integers forming a magic square with the magic constant 3c so long as 0 < a < b < c − a and b ≠ 2a. Moreover, every 3×3 magic square of distinct positive integers is of this form.
In 1997 Lee Sallows discovered that leaving aside rotations and reflections, then every distinct parallelogram drawn on the Argand diagram defines a unique 3×3 magic square, and vice versa, a result that had never previously been noted.
A method for constructing a magic square of odd order
A method for constructing magic squares of odd order was published by the French diplomat de la Loubère in his book, A new historical relation of the kingdom of Siam (Du Royaume de Siam, 1693), in the chapter entitled The problem of the magical square according to the Indians. The method operates as follows:
The method prescribes starting in the central column of the first row with the number 1. After that, the fundamental movement for filling the squares is diagonally up and right, one step at a time. If a square is filled with a multiple of the order n, one moves vertically down one square instead, then continues as before. When an "up and to the right" move would leave the square, it is wrapped around to the last row or first column, respectively.
Starting from other squares rather than the central column of the first row is possible, but then only the row and column sums will be identical and result in a magic sum, whereas the diagonal sums will differ. The result will thus be a semimagic square and not a true magic square. Moving in directions other than north east can also result in magic squares.
A method of constructing a magic square of doubly even order
Doubly even means that n is an even multiple of an even integer; or 4p (e.g. 4, 8, 12), where p is an integer.
Generic pattern
All the numbers are written in order from left to right across each row in turn, starting from the top left hand corner. Numbers are then either retained in the same place or interchanged with their diametrically opposite numbers in a certain regular pattern. In the magic square of order four, the numbers in the four central squares and one square at each corner are retained in the same place and the others are interchanged with their diametrically opposite numbers.
A construction of a magic square of order 4
Starting from top left, go left to right through each row of the square, counting each cell from 1 to 16 and filling the cells along the diagonals with its corresponding number. Once the bottom right cell is reached, continue by going right to left, starting from the bottom right of the table through each row, and fill in the non-diagonal cells counting up from 1 to 16 with its corresponding number. As shown below:
An extension of the above example for Orders 8 and 12
First generate a pattern table, where a '1' indicates selecting from the square where the numbers are written in order 1 to n2 (left-to-right, top-to-bottom), and a '0' indicates selecting from the square where the numbers are written in reverse order n2 to 1. For M = 4, the pattern table is as shown below (third matrix from left). With the unaltered cells (cells with '1') shaded, a criss-cross pattern is obtained.
The patterns are a) there are equal number of '1's and '0's in each row and column; b) each row and each column are "palindromic"; c) the left- and right-halves are mirror images; and d) the top- and bottom-halves are mirror images (c and d imply b). The pattern table can be denoted using hexadecimals as (9, 6, 6, 9) for simplicity (1-nibble per row, 4 rows). The simplest method of generating the required pattern for higher ordered doubly even squares is to copy the generic pattern for the fourth-order square in each four-by-four sub-squares.
For M = 8, possible choices for the pattern are (99, 66, 66, 99, 99, 66, 66, 99); (3C, 3C, C3, C3, C3, C3, 3C, 3C); (A5, 5A, A5, 5A, 5A, A5, 5A, A5) (2-nibbles per row, 8 rows).
For M = 12, the pattern table (E07, E07, E07, 1F8, 1F8, 1F8, 1F8, 1F8, 1F8, E07, E07, E07) yields a magic square (3-nibbles per row, 12 rows.) It is possible to count the number of choices one has based on the pattern table, taking rotational symmetries into account.
Method of superposition
The earliest discovery of the superposition method was made by the Indian mathematician Narayana in the 14th century. The same method was later re-discovered and studied in early 18th century Europe by de la Loubere, Poignard, de La Hire, and Sauveur; and the method is usually referred to as de la Hire's method. Although Euler's work on magic square was unoriginal, he famously conjectured the impossibility of constructing the evenly odd ordered mutually orthogonal Graeco-Latin squares. This conjecture was disproved in the mid 20th century. For clarity of exposition, two important variations of this method can be distinguished.
Euler's method
This method consists in constructing two preliminary squares, which when added together gives the magic square. As a running example, a 3×3 magic square is considered. Each number of the 3×3 natural square by a pair of numbers can be labeled as
where every pair of Greek and Latin alphabets, e.g. αa, are meant to be added together, i.e. αa = α + a. Here, (α, β, γ) = (0, 3, 6) and (a, b, c) = (1, 2, 3). The numbers 0, 3, and 6 are referred to as the root numbers while the numbers 1, 2, and 3 are referred to as the primary numbers. An important general constraint here is
a Greek letter is paired with a Latin letter only once.
Thus, the original square can now be split into two simpler squares:
The lettered squares are referred to as Greek square or Latin square if they are filled with Greek or Latin letters, respectively. A magic square can be constructed by ensuring that the Greek and Latin squares are magic squares too. The converse of this statement is also often, but not always (e.g. bordered magic squares), true: A magic square can be decomposed into a Greek and a Latin square, which are themselves magic squares. Thus the method is useful for both synthesis as well as analysis of a magic square. Lastly, by examining the pattern in which the numbers are laid out in the finished square, it is often possible to come up with a faster algorithm to construct higher order squares that replicate the given pattern, without the necessity of creating the preliminary Greek and Latin squares.
During the construction of the 3×3 magic square, the Greek and Latin squares with just three unique terms are much easier to deal with than the original square with nine different terms. The row sum and the column sum of the Greek square will be the same, α + β + γ, if
each letter appears exactly once in a given column or a row.
This can be achieved by cyclic permutation of α, β, and γ. Satisfaction of these two conditions ensures that the resulting square is a semi-magic square; and such Greek and Latin squares are said to be mutually orthogonal to each other. For a given order n, there are at most n - 1 squares in a set of mutually orthogonal squares, not counting the variations due to permutation of the symbols. This upper bound is exact when n is a prime number.
In order to construct a magic square, we should also ensure that the diagonals sum to magic constant. For this, we have a third condition:
either all the letters should appear exactly once in both the diagonals; or in case of odd ordered squares, one of the diagonals should consist entirely of the middle term, while the other diagonal should have all the letters exactly once.
The mutually orthogonal Greek and Latin squares that satisfy the first part of the third condition (that all letters appear in both the diagonals) are said to be mutually orthogonal doubly diagonal Graeco-Latin squares.
Odd squares: For the 3×3 odd square, since α, β, and γ are in arithmetic progression, their sum is equal to the product of the square's order and the middle term, i.e. α + β + γ = 3 β. Thus, the diagonal sums will be equal if we have βs in the main diagonal and α, β, γ in the skew diagonal. Similarly, for the Latin square. The resulting Greek and Latin squares and their combination will be as below. The Latin square is just a 90 degree anti-clockwise rotation of the Greek square (or equivalently, flipping about the vertical axis) with the corresponding letters interchanged. Substituting the values of the Greek and Latin letters will give the 3×3 magic square.
For the odd squares, this method explains why the Siamese method (method of De la Loubere) and its variants work. This basic method can be used to construct odd ordered magic squares of higher orders. To summarise:
For odd ordered squares, to construct Greek square, place the middle term along the main diagonal, and place the rest of the terms along the skew diagonal. The remaining empty cells are filled by diagonal moves. The Latin square can be constructed by rotating or flipping the Greek square, and replacing the corresponding alphabets. The magic square is obtained by adding the Greek and Latin squares.
A peculiarity of the construction method given above for the odd magic squares is that the middle number (n2 + 1)/2 will always appear at the center cell of the magic square. Since there are (n - 1)! ways to arrange the skew diagonal terms, we can obtain (n - 1)! Greek squares this way; same with the Latin squares. Also, since each Greek square can be paired with (n - 1)! Latin squares, and since for each of Greek square the middle term may be arbitrarily placed in the main diagonal or the skew diagonal (and correspondingly along the skew diagonal or the main diagonal for the Latin squares), we can construct a total of 2 × (n - 1)! × (n - 1)! magic squares using this method. For n = 3, 5, and 7, this will give 8, 1152, and 1,036,800 different magic squares, respectively. Dividing by 8 to neglect equivalent squares due to rotation and reflections, we obtain 1, 144, and 129,600 essentially different magic squares, respectively.
As another example, the construction of 5×5 magic square is given. Numbers are directly written in place of alphabets. The numbered squares are referred to as primary square or root square if they are filled with primary numbers or root numbers, respectively. The numbers are placed about the skew diagonal in the root square such that the middle column of the resulting root square has 0, 5, 10, 15, 20 (from bottom to top). The primary square is obtained by rotating the root square counter-clockwise by 90 degrees, and replacing the numbers. The resulting square is an associative magic square, in which every pair of numbers symmetrically opposite to the center sum up to the same value, 26. For e.g., 16+10, 3+23, 6+20, etc. In the finished square, 1 is placed at center cell of bottom row, and successive numbers are placed via elongated knight's move (two cells right, two cells down), or equivalently, bishop's move (two cells diagonally down right). When a collision occurs, the break move is to move one cell up. All the odd numbers occur inside the central diamond formed by 1, 5, 25 and 21, while the even numbers are placed at the corners. The occurrence of the even numbers can be deduced by copying the square to the adjacent sides. The even numbers from four adjacent squares will form a cross.
A variation of the above example, where the skew diagonal sequence is taken in different order, is given below. The resulting magic square is the flipped version of the famous Agrippa's Mars magic square. It is an associative magic square and is the same as that produced by Moschopoulos's method. Here the resulting square starts with 1 placed in the cell which is to the right of the centre cell, and proceeds as De la Loubere's method, with downwards-right move. When a collision occurs, the break move is to shift two cells to the right.
In the previous examples, for the Greek square, the second row can be obtained from the first row by circularly shifting it to the right by one cell. Similarly, the third row is a circularly shifted version of the second row by one cell to the right; and so on. Likewise, the rows of the Latin square is circularly shifted to the left by one cell. The row shifts for the Greek and Latin squares are in mutually opposite direction. It is possible to circularly shift the rows by more than one cell to create the Greek and Latin square.
For odd ordered squares, whose order is not divisible by three, we can create the Greek squares by shifting a row by two places to the left or to the right to form the next row. The Latin square is made by flipping the Greek square along the main diagonal and interchanging the corresponding letters. This gives us a Latin square whose rows are created by shifting the row in the direction opposite to that of the Greek square. A Greek square and a Latin square should be paired such that their row shifts are in mutually opposite direction. The magic square is obtained by adding the Greek and Latin squares. When the order also happens to be a prime number, this method always creates pandiagonal magic square.
This essentially re-creates the knight's move. All the letters will appear in both the diagonals, ensuring correct diagonal sum. Since there are n! permutations of the Greek letters by which we can create the first row of the Greek square, there are thus n! Greek squares that can be created by shifting the first row in one direction. Likewise, there are n! such Latin squares created by shifting the first row in the opposite direction. Since a Greek square can be combined with any Latin square with opposite row shifts, there are n! × n! such combinations. Lastly, since the Greek square can be created by shifting the rows either to the left or to the right, there are a total of 2 × n! × n! magic squares that can be formed by this method. For n = 5 and 7, since they are prime numbers, this method creates 28,800 and 50,803,200 pandiagonal magic squares. Dividing by 8 to neglect equivalent squares due to rotation and reflections, we obtain 3,600 and 6,350,400 equivalent squares. Further dividing by n2 to neglect equivalent panmagic squares due to cyclic shifting of rows or columns, we obtain 144 and 129,600 essentially different panmagic squares. For order 5 squares, these are the only panmagic square there are. The condition that the square's order not be divisible by 3 means that we cannot construct squares of orders 9, 15, 21, 27, and so on, by this method.
In the example below, the square has been constructed such that 1 is at the center cell. In the finished square, the numbers can be continuously enumerated by the knight's move (two cells up, one cell right). When collision occurs, the break move is to move one cell up, one cell left. The resulting square is a pandiagonal magic square. This square also has a further diabolical property that any five cells in quincunx pattern formed by any odd sub-square, including wrap around, sum to the magic constant, 65. For e.g., 13+7+1+20+24, 23+1+9+15+17, 13+21+10+19+2 etc. Also the four corners of any 5×5 square and the central cell, as well as the middle cells of each side together with the central cell, including wrap around, give the magic sum: 13+10+19+22+1 and 20+24+12+8+1. Lastly the four rhomboids that form elongated crosses also give the magic sum: 23+1+9+24+8, 15+1+17+20+12, 14+1+18+13+19, 7+1+25+22+10. Such squares with 1 at the center cell are also called God's magic squares in Islamic amulet design, where the center cell is either left blank or filled with God's name.
We can also combine the Greek and Latin squares constructed by different methods. In the example below, the primary square is made using knight's move. We have re-created the magic square obtained by De la Loubere's method. As before, we can form 8 × (n - 1)! × n! magic squares by this combination. For n = 5 and 7, this will create 23,040 and 29,030,400 magic squares. After dividing by 8 in order to neglect equivalent squares due to rotation and reflection, we get 2,880 and 3,628,800 squares.
For order 5 squares, these three methods give a complete census of the number of magic squares that can be constructed by the method of superposition. Neglecting the rotation and reflections, the total number of magic squares of order 5 produced by the superposition method is 144 + 3,600 + 2,880 = 6,624.
Even squares: We can also construct even ordered squares in this fashion. Since there is no middle term among the Greek and Latin alphabets for even ordered squares, in addition to the first two constraint, for the diagonal sums to yield the magic constant, all the letters in the alphabet should appear in the main diagonal and in the skew diagonal.
An example of a 4×4 square is given below. For the given diagonal and skew diagonal in the Greek square, the rest of the cells can be filled using the condition that each letter appear only once in a row and a column.
Using these two Graeco-Latin squares, we can construct 2 × 4! × 4! = 1,152 magic squares. Dividing by 8 to eliminate equivalent squares due to rotation and reflections, we get 144 essentially different magic squares of order 4. These are the only magic squares constructible by the Euler method, since there are only two mutually orthogonal doubly diagonal Graeco-Latin squares of order 4.
Similarly, an 8×8 magic square can be constructed as below. Here the order of appearance of the numbers is not important; however the quadrants imitate the layout pattern of the 4×4 Graeco-Latin squares.
Euler's method has given rise to the study of Graeco-Latin squares. Euler's method for constructing magic squares is valid for any order except 2 and 6.
Variations: Magic squares constructed from mutually orthogonal doubly diagonal Graeco-Latin squares are interesting in themselves since the magic property emerges from the relative position of the alphabets in the square, and not due to any arithmetic property of the value assigned to them. This means that we can assign any value to the alphabets of such squares and still obtain a magic square. This is the basis for constructing squares that display some information (e.g. birthdays, years, etc.) in the square and for creating "reversible squares". For example, we can display the number π ≈ at the bottom row of a 4×4 magic square using the Graeco-Latin square given above by assigning (α, β, γ, δ) = (10, 0, 90, 15) and (a, b, c, d) = (0, 2, 3, 4). We will obtain the following non-normal magic square with the magic sum 124:
Narayana-De la Hire's method for even orders
Narayana-De la Hire's method for odd square is the same as that of Euler's. However, for even squares, we drop the second requirement that each Greek and Latin letter appear only once in a given row or column. This allows us to take advantage of the fact that the sum of an arithmetic progression with an even number of terms is equal to the sum of two opposite symmetric terms multiplied by half the total number of terms. Thus, when constructing the Greek or Latin squares,
for even ordered squares, a letter can appear n/2 times in a column but only once in a row, or vice versa.
As a running example, if we take a 4×4 square, where the Greek and Latin terms have the values (α, β, γ, δ) = (0, 4, 8, 12) and (a, b, c, d) = (1, 2, 3, 4), respectively, then we have α + β + γ + δ = 2 (α + δ) = 2 (β + γ). Similarly, a + b + c + d = 2 (a + d) = 2 (b + c). This means that the complementary pair α and δ (or β and γ) can appear twice in a column (or a row) and still give the desired magic sum. Thus, we can construct:
For even ordered squares, the Greek magic square is made by first placing the Greek alphabets along the main diagonal in some order. The skew diagonal is then filled in the same order or by picking the terms that are complementary to the terms in the main diagonal. Finally, the remaining cells are filled column wise. Given a column, we use the complementary terms in the diagonal cells intersected by that column, making sure that they appear only once in a given row but n/2 times in the given column. The Latin square is obtained by flipping or rotating the Greek square and interchanging the corresponding alphabets. The final magic square is obtained by adding the Greek and Latin squares.
In the example given below, the main diagonal (from top left to bottom right) is filled with sequence ordered as α, β, γ, δ, while the skew diagonal (from bottom left to top right) filled in the same order. The remaining cells are then filled column wise such that the complementary letters appears only once within a row, but twice within a column. In the first column, since α appears on the 1st and 4th row, the remaining cells are filled with its complementary term δ. Similarly, the empty cells in the 2nd column are filled with γ; in 3rd column β; and 4th column α. Each Greek letter appears only once along the rows, but twice along the columns. As such, the row sums are α + β + γ + δ while the column sums are either 2 (α + δ) or 2 (β + γ). Likewise for the Latin square, which is obtained by flipping the Greek square along the main diagonal and interchanging the corresponding letters.
The above example explains why the "criss-cross" method for doubly even magic square works. Another possible 4×4 magic square, which is also pan-diagonal as well as most-perfect, is constructed below using the same rule. However, the diagonal sequence is chosen such that all four letters α, β, γ, δ appear inside the central 2×2 sub-square. Remaining cells are filled column wise such that each letter appears only once within a row. In the 1st column, the empty cells need to be filled with one of the letters selected from the complementary pair α and δ. Given the 1st column, the entry in the 2nd row can only be δ since α is already there in the 2nd row; while, in the 3rd row the entry can only be α since δ is already present in the 3rd row. We proceed similarly until all cells are filled. The Latin square given below has been obtained by flipping the Greek square along the main diagonal and replacing the Greek alphabets with corresponding Latin alphabets.
We can use this approach to construct singly even magic squares as well. However, we have to be more careful in this case since the criteria of pairing the Greek and Latin alphabets uniquely is not automatically satisfied. Violation of this condition leads to some missing numbers in the final square, while duplicating others. Thus, here is an important proviso:
For singly even squares, in the Greek square, check the cells of the columns which is vertically paired to its complement. In such a case, the corresponding cell of the Latin square must contain the same letter as its horizontally paired cell.
Below is a construction of a 6×6 magic square, where the numbers are directly given, rather than the alphabets. The second square is constructed by flipping the first square along the main diagonal. Here in the first column of the root square the 3rd cell is paired with its complement in the 4th cells. Thus, in the primary square, the numbers in the 1st and 6th cell of the 3rd row are same. Likewise, with other columns and rows. In this example the flipped version of the root square satisfies this proviso.
As another example of a 6×6 magic square constructed this way is given below. Here the diagonal entries are arranged differently. The primary square is constructed by flipping the root square about the main diagonal. In the second square the proviso for singly even square is not satisfied, leading to a non-normal magic square (third square) where the numbers 3, 13, 24, and 34 are duplicated while missing the numbers 4, 18, 19, and 33.
The last condition is a bit arbitrary and may not always need to be invoked, as in this example, where in the root square each cell is vertically paired with its complement:
As one more example, we have generated an 8×8 magic square. Unlike the criss-cross pattern of the earlier section for evenly even square, here we have a checkered pattern for the altered and unaltered cells. Also, in each quadrant the odd and even numbers appear in alternating columns.
Variations: A number of variations of the basic idea are possible: a complementary pair can appear n/2 times or less in a column. That is, a column of a Greek square can be constructed using more than one complementary pair. This method allows us to imbue the magic square with far richer properties. The idea can also be extended to the diagonals too. An example of an 8×8 magic square is given below. In the finished square each of four quadrants are pan-magic squares as well, each quadrant with same magic constant 130.
Method of borders
Bordering method for order 3
In this method, the objective is to wrap a border around a smaller magic square which serves as a core. Consider the 3×3 square for example. Subtracting the middle number 5 from each number 1, 2, ..., 9, we obtain 0, ± 1, ± 2, ± 3, and ± 4, which we will, for lack of better words, following S. Harry White, refer to as bone numbers. The magic constant of a magic square, which we will refer to as the skeleton square, made by these bone numbers will be zero since adding all the rows of a magic square will give nM = Σ k = 0; thus M = 0.
It is not difficult to argue that the middle number should be placed at the center cell: let x be the number placed in the middle cell, then the sum of the middle column, middle row, and the two diagonals give Σ k + 3 x = 4 M. Since Σ k = 3 M, we have x = M / 3. Here M = 0, so x = 0.
Putting the middle number 0 in the center cell, we want to construct a border such that the resulting square is magic. Let the border be given by:
Since the sum of each row, column, and diagonals must be a constant (which is zero), we have
a + a* = 0,
b + b* = 0,
u + u* = 0,
v + v* = 0.
Now, if we have chosen a, b, u, and v, then we have a* = - a, b* = - b, u* = - u, and v* = - v. This means that if we assign a given number to a variable, say a = 1, then its complement will be assigned to a*, i.e. a* = - 1. Thus out of eight unknown variables, it is sufficient to specify the value of only four variables. We will consider a, b, u, and v as independent variables, while a*, b*, u*, and v* as dependent variables. This allows us to consider a bone number ± x as a single number regardless of sign because (1) its assignment to a given variable, say a, will automatically imply that the same number of opposite sign will be shared with its complement a*, and (2) two independent variables, say a and b, cannot be assigned the same bone number. But how should we choose a, b, u, and v? We have the sum of the top row and the sum of the right column as
u + a + v = 0,
v + b + u* = 0.
Since 0 is an even number, there are only two ways that the sum of three integers will yield an even number: 1) if all three were even, or 2) if two were odd and one was even. Since in our choice of numbers we only have two even non-zero number (± 2 and ± 4), the first statement is false. Hence, it must be the case that the second statement is true: that two of the numbers are odd and one even.
The only way that both the above two equations can satisfy this parity condition simultaneously, and still be consistent with the set of numbers we have, is when u and v are odd. For on the contrary, if we had assumed u and a to be odd and v to be even in the first equation, then u* = - u will be odd in the second equation, making b odd as well, in order to satisfy the parity condition. But this requires three odd numbers (u, a, and b), contradicting the fact that we only have two odd numbers (± 1 and ± 3) which we can use. This proves that the odd bone numbers occupy the corners cells. When converted to normal numbers by adding 5, this implies that the corners of a 3×3 magic square are all occupied by even numbers.
Thus, taking u = 1 and v = 3, we have a = - 4 and b = - 2. Hence, the finished skeleton square will be as in the left. Adding 5 to each number, we get the finished magic square.
Similar argument can be used to construct larger squares. Since there does not exist a 2×2 magic square around which we can wrap a border to construct a 4×4 magic square, the next smallest order for which we can construct bordered square is the order 5.
Bordering method for order 5
Consider the fifth-order square. For this, we have a 3×3 magic core, around which we will wrap a magic border. The bone numbers to be used will be ± 5, ± 6, ± 7, ± 8, ± 9, ± 10, ± 11, and ± 12. Disregarding the signs, we have 8 bone numbers, 4 of which are even and 4 of which are odd. In general, for a square of any order n, there will be 4(n - 1) border cells, which are to be filled using 2(n - 1) bone numbers. Let the magic border be given as
As before, we should
place a bone number and its complement opposite to each other, so that the magic sum will be zero.
It is sufficient to determine the numbers u, v, a, b, c, d, e, f to describe the magic border. As before, we have the two constraint equations for the top row and right column:
u + a + b + c + v = 0
v + d + e + f + u* = 0.
Multiple solutions are possible. The standard procedure is to
first try to determine the corner cells, after which we will try to determine the rest of the border.
There are 28 ways of choosing two numbers from the set of 8 bone numbers for the corner cells u and v. However, not all pairs are admissible. Among the 28 pairs, 16 pairs are made of an even and an odd number, 6 pairs have both as even numbers, while 6 pairs have them both as odd numbers.
We can prove that the corner cells u and v cannot have an even and an odd number. This is because if this were so, then the sums u + v and v + u* will be odd, and since 0 is an even number, the sums a + b + c and d + e + f should be odd as well. The only way that the sum of three integers will result in an odd number is when 1) two of them are even and one is odd, or 2) when all three are odd. Since the corner cells are assumed to be odd and even, neither of these two statements are compatible with the fact that we only have 3 even and 3 odd bone numbers at our disposal. This proves that u and v cannot have different parity. This eliminates 16 possibilities.
Using similar type reasoning we can also draw some conclusions about the sets {a, b, c} and {d, e, f}. If u and v are both even, then both the sets should have two odd numbers and one even number. If u and v are both odd, then one of the sets should have three even numbers while the other set should have one even number and two odd numbers.
As a running example, consider the case when both u and v are even. The 6 possible pairs are: (6, 8), (6, 10), (6, 12), (8, 10), (8, 12), and (10, 12). Since the sums u + v and v + u* are even, the sums a + b + c and d + e + f should be even as well. The only way that the sum of three integers will result in an even number is when 1) two of them are odd and one is even, or 2) when all three are even. The fact that the two corner cells are even means that we have only 2 even numbers at our disposal. Thus, the second statement is not compatible with this fact. Hence, it must be the case that the first statement is true: two of the three numbers should be odd, while one be even.
Now let a, b, d, e be odd numbers while c and f be even numbers. Given the odd bone numbers at our disposal: ± 5, ± 7, ± 9, and ± 11, their differences range from D = { ± 2, ± 4, ± 6} while their sums range from S = {± 12, ± 14, ± 16, ± 18, ± 20}. It is also useful to have a table of their sum and differences for later reference. Now, given the corner cells (u, v), we can check its admissibility by checking if the sums u + v + c and v + u* + f fall within the set D or S. The admissibility of the corner numbers is a necessary but not a sufficient condition for the solution to exist.
For example, if we consider the pair (u, v) = (8, 12), then u + v = 20 and v + u* = 6; and we will have ± 6 and ± 10 even bone numbers at our disposal. Taking c = ± 6, we have the sum u + v + c to be 26 and 14, depending on the sign of ± 6 taken, both of which do not fall within the sets D or S. Likewise, taking c = ± 10, we have the sum u + v + c to be 30 and 10, both of which again do not fall within the sets D or S. Thus, the pair (8, 12) is not admissible. By similar process of reasoning, we can also rule out the pair (6, 12).
As another example, if we consider the pair (u, v) = (10, 12), then u + v = 22 and v + u* = 2; and we will have ± 6 and ± 8 even bone numbers at our disposal. Taking c = ± 6, we have the sum u + v + c to be 28 and 16. While 28 does not fall within the sets D or S, 16 falls in set S. By inspection, we find that if (a, b) = (-7, -9), then a + b = -16; and it will satisfy the first constraint equation. Also, taking f = ± 8, we have the sum v + u* + f to be 10 and -6. While 10 does not fall within the sets D or S, -6 falls in set D. Since -7 and -9 have already been assigned to a and b, clearly (d, e) = (-5, 11) so that d + e = 6; and it will satisfy the second constraint equation.
Likewise, taking c = ± 8, we have the sum u + v + c to be 30 and 14. While 30 does not fall within the sets D or S, 14 falls in set S. By inspection, we find that if (a, b) = (-5, -9), then a + b = -14. Also, taking f = ± 6, we have the sum v + u* + f to be 8 and -4. While 8 does not fall within the sets D or S, -4 falls in set D. Clearly, (d, e) = (-7, 11) so that d + e = 4, and the second constraint equation will be satisfied.
Hence the corner pair (u, v) = (10, 12) is admissible; and it admits two solutions: (a, b, c, d, e, f) = (-7, -9, -6, -5, 11, -8) and (a, b, c, d, e, f) = ( -5, -9, -8, -7, 11, -6). The finished skeleton squares are given below. The magic square is obtained by adding 13 to each cells.
Using similar process of reasoning, we can construct the following table for the values of u, v, a, b, c, d, e, f expressed as bone numbers as given below. There are only 6 possible choices for the corner cells, which leads to 10 possible border solutions.
Given this group of 10 borders, we can construct 10×8×(3!)2 = 2880 essentially different bordered magic squares. Here the bone numbers ± 5, ..., ± 12 were consecutive. More bordered squares can be constructed if the numbers are not consecutive. If non-consecutive bone numbers were also used, then there are a total of 605 magic borders. Thus, the total number of order 5 essentially different bordered magic squares (with consecutive and non-consecutive numbers) is 174,240. See history. The number of fifth-order magic squares constructible via the bordering method is about 26 times larger than via the superposition method.
Continuous enumeration methods
Exhaustive enumeration of all the borders of a magic square of a given order, as done previously, is very tedious. As such a structured solution is often desirable, which allows us to construct a border for a square of any order. Below we give three algorithms for constructing border for odd, doubly even, and singly even squares. These continuous enumeration algorithms were discovered in 10th century by Arab scholars; and their earliest surviving exposition comes from the two treatises by al-Buzjani and al-Antaki, although they themselves were not the discoverers. Since then many more such algorithms have been discovered.
Odd-ordered squares: The following is the algorithm given by al-Buzjani to construct a border for odd squares. A peculiarity of this method is that for order n square, the two adjacent corners are numbers n - 1 and n + 1.
Starting from the cell above the lower left corner, we put the numbers alternately in left column and bottom row until we arrive at the middle cell. The next number is written in the middle cell of the bottom row just reached, after which we fill the cell in the upper left corner, then the middle cell of the right column, then the upper right corner. After this, starting from the cell above middle cell of the right column already filled, we resume the alternate placement of the numbers in the right column and the top row. Once half of the border cells are filled, the other half are filled by numbers complementary to opposite cells. The subsequent inner borders is filled in the same manner, until the square of order 3 is filled.
Below is an example for 9th-order square.
Doubly even order: The following is the method given by al-Antaki. Consider an empty border of order n = 4k with k ≥ 3. The peculiarity of this algorithm is that the adjacent corner cells are occupied by numbers n and n - 1.
Starting at the upper left corner cell, we put the successive numbers by groups of four, the first one next to the corner, the second and the third on the bottom, and the fourth at the top, and so on until there remains in the top row (excluding the corners) six empty cells. We then write the next two numbers above and the next four below. We then fill the upper corners, first left then right. We place the next number below the upper right corner in the right column, the next number on the other side in the left column. We then resume placing groups of four consecutive numbers in the two columns as before. Once half of the border cells are filled, the other half are filled by numbers complementary to opposite cells.
The example below gives the border for order 16 square.
For order 8 square, we just begin directly with the six cells.
Singly even order: For singly even order, we have the algorithm given by al-Antaki. Here the corner cells are occupied by n and n - 1. Below is an example of 10th-order square.
Start by placing 1 at the bottom row next to the left corner cell, then place 2 in the top row. After this, place 3 at the bottom row and turn around the border in anti-clockwise direction placing the next numbers, until n - 2 is reached on the right column. The next two numbers are placed in the upper corners (n - 1 in upper left corner and n in upper right corner). Then, the next two numbers are placed on the left column, then we resume the cyclic placement of the numbers until half of all the border cells are filled. Once half of the border cells are filled, the other half are filled by numbers complementary to opposite cells.
Method of composition
For squares of order m × n where m, n > 2
This is a method reminiscent of the Kronecker product of two matrices, that builds an nm × nm magic square from an n × n magic square and an m × m magic square. The "product" of two magic squares creates a magic square of higher order than the two multiplicands. Let the two magic squares be of orders m and n. The final square will be of order m × n. Divide the square of order m × n into m × m sub-squares, such that there are a total of n2 such sub-squares. In the square of order n, reduce by 1 the value of all the numbers. Multiply these reduced values by m2, and place the results in the corresponding sub-squares of the m × n whole square. The squares of order m are added n2 times to the sub-squares of the final square. The peculiarity of this construction method is that each magic subsquare will have different magic sums. The square made of such magic sums from each magic subsquare will again be a magic square. The smallest composite magic square of order 9, composed of two order 3 squares is given below.
Since each of the 3×3 sub-squares can be independently rotated and reflected into 8 different squares, from this single 9×9 composite square we can derive 89 = 134,217,728 essentially different 9×9 composite squares. Plenty more composite magic squares can also be derived if we select non-consecutive numbers in the magic sub-squares, like in Yang Hui's version of the 9×9 composite magic square. The next smallest composite magic squares of order 12, composed of magic squares of order 3 and 4 are given below.
For the base squares, there is only one essentially different 3rd order square, while there 880 essentially different 4th-order squares that we can choose from. Each pairing can produce two different composite squares. Since each magic sub-squares in each composite square can be expressed in 8 different forms due to rotations and reflections, there can be 1×880×89 + 880×1×816 ≈ 2.476×1017 essentially different 12×12 composite magic squares created this way, with consecutive numbers in each sub-square. In general, if there are cm and cn essentially different magic squares of order m and n, then we can form cm × cn × ( 8m2 + 8n2) composite squares of order mn, provided m ≠ n. If m = n, then we can form (cm)2 × 8m2 composite squares of order m2.
For squares of doubly even order
When the squares are of doubly even order, we can construct a composite magic square in a manner more elegant than the above process, in the sense that every magic subsquare will have the same magic constant. Let n be the order of the main square and m the order of the equal subsquares. The subsquares are filled one by one, in any order, with a continuous sequence of m2/2 smaller numbers (i.e. numbers less than or equal to n2/2) together with their complements to n2 + 1. Each subsquare as a whole will yield the same magic sum. The advantage of this type of composite square is that each subsquare is filled in the same way and their arrangement is arbitrary. Thus, the knowledge of a single construction of even order will suffice to fill the whole square. Furthermore, if the subsquares are filled in the natural sequence, then the resulting square will be pandiagonal. The magic sum of the subsquares is related to the magic sum of the whole square by where n = km.
In the examples below, we have divided the order 12 square into nine subsquares of order 4 filled each with eight smaller numbers and, in the corresponding bishop's cells (two cells diagonally across, including wrap arounds, in the 4×4 subsquare), their complements to n2 + 1 = 145. Each subsquare is pandiagonal with magic constant 290; while the whole square on the left is also pandiagonal with magic constant 870.
In another example below, we have divided the order 12 square into four order 6 squares. Each of the order 6 squares are filled with eighteen small numbers and their complements using bordering technique given by al-Antaki. If we remove the shaded borders of the order 6 subsquares and form an order 8 square, then this order 8 square is again a magic square. In its full generality, we can take any m2/2 smaller numbers together with their complements to n2 + 1 to fill the subsquares, not necessarily in continuous sequence.
Medjig-method for squares of even order 2n, where n > 2
In this method a magic square is "multiplied" with a medjig square to create a larger magic square. The namesake of this method derives from mathematical game called medjig created by Willem Barink in 2006, although the method itself is much older. An early instance of a magic square constructed using this method occurs in Yang Hui's text for order 6 magic square. The LUX method to construct singly even magic squares is a special case of the medjig method, where only 3 out of 24 patterns are used to construct the medjig square.
The pieces of the medjig puzzle are 2×2 squares on which the numbers 0, 1, 2 and 3 are placed. There are three basic patterns by which the numbers 0, 1, 2 and 3 can be placed in a 2×2 square, where 0 is at the top left corner:
Each pattern can be reflected and rotated to obtain 8 equivalent patterns, giving us a total of 3×8 = 24 patterns. The aim of the puzzle is to take n2 medjig pieces and arrange them in an n × n medjig square in such a way that each row, column, along with the two long diagonals, formed by the medjig square sums to 3n, the magic constant of the medjig square. An n × n medjig square can create a 2n × 2n magic square where n > 2.
Given an n×n medjig square and an n×n magic square base, a magic square of order 2n×2n can be constructed as follows:
Each cell of an n×n magic square is associated with a corresponding 2×2 subsquare of the medjig square
Fill each 2×2 subsquares of the medjig square with the four numbers from 1 to 4n2 that equal the original number modulo n2, i.e. x+n2y where x is the corresponding number from the magic square and y is a number from 0 to 3 in the 2×2 subsquares.
Assuming that we have an initial magic square base, the challenge lies in constructing a medjig square. For reference, the sums of each medjig piece along the rows, columns and diagonals, denoted in italics, are:
Doubly even squares: The smallest even ordered medjig square is of order 2 with magic constant 6. While it is possible to construct a 2×2 medjig square, we cannot construct a 4×4 magic square from it since 2×2 magic squares required to "multiply" it does not exist. Nevertheless, it is worth constructing these 2×2 medjig squares. The magic constant 6 can be partitioned into two parts in three ways as 6 = 5 + 1 = 4 + 2 = 3 + 3. There exist 96 such 2×2 medjig squares. In the examples below, each 2×2 medjig square is made by combining different orientations of a single medjig piece.
We can use the 2×2 medjig squares to construct larger even ordered medjig squares. One possible approach is to simply combine the 2×2 medjig squares together. Another possibility is to wrap a smaller medjig square core with a medjig border. The pieces of a 2×2 medjig square can form the corner pieces of the border. Yet another possibility is to append a row and a column to an odd ordered medjig square. An example of an 8×8 magic square is constructed below by combining four copies of the left most 2×2 medjig square given above:
The next example is constructed by bordering a 2×2 medjig square core.
Singly even squares: Medjig square of order 1 does not exist. As such, the smallest odd ordered medjig square is of order 3, with magic constant 9. There are only 7 ways of partitioning the integer 9, our magic constant, into three parts. If these three parts correspond to three of the medjig pieces in a row, column or diagonal, then the relevant partitions for us are:
9 = 1 + 3 + 5 = 1 + 4 + 4 = 2 + 3 + 4 = 2 + 2 + 5 = 3 + 3 + 3.
A 3×3 medjig square can be constructed with some trial-and-error, as in the left most square below. Another approach is to add a row and a column to a 2×2 medjig square. In the middle square below, a left column and bottom row has been added, creating an L-shaped medjig border, to a 2×2 medjig square given previously. The right most square below is essentially same as the middle square, except that the row and column has been added in the middle to form a cross while the pieces of 2×2 medjig square are placed at the corners.
Once a 3×3 medjig square has been constructed, it can be converted into a 6×6 magic square. For example, using the left most 3×3 medjig square given above:
There are 1,740,800 such 3×3 medjig squares. An easy approach to construct higher order odd medjig square is by wrapping a smaller odd ordered medjig square with a medjig border, just as with even ordered medjig squares. Another approach is to append a row and a column to an even ordered medjig square. Approaches such as the LUX method can also be used. In the example below, a 5×5 medjig square is created by wrapping a medjig border around a 3×3 medjig square given previously:
Solving partially completed magic squares
Solving partially completed magic squares is a popular mathematical pastime. The techniques needed are similar to those used in Sudoku or KenKen puzzles, and involve deducing the values of unfilled squares using logic and permutation group theory (Sudoku grids are not magic squares but are based on a related idea called Graeco-Latin squares).
Variations of the magic square
Extra constraints
Certain extra restrictions can be imposed on magic squares.
If raising each number to the nth power yields another magic square, the result is a bimagic (n = 2), a trimagic (n = 3), or, in general, a multimagic square.
A magic square in which the number of letters in the name of each number in the square generates another magic square is called an alphamagic square.
There are magic squares consisting entirely of primes. Rudolf Ondrejka (1928–2001) discovered the following 3×3 magic square of primes, in this case nine Chen primes:
The Green–Tao theorem implies that there are arbitrarily large magic squares consisting of primes.
The following "reversible magic square" has a magic constant of 264 both upside down and right way up:
When the extra constraint is to display some date, especially a birth date, then such magic squares are called birthday magic square. An early instance of such birthday magic square was created by Srinivasa Ramanujan. He created a 4×4 square in which he entered his date of birth in D–M–C-Y format in the top row and the magic happened with additions and subtractions of numbers in squares. Not only do the rows, columns, and diagonals add up to the same number, but the four corners, the four middle squares (17, 9, 24, 89), the first and last rows two middle numbers (12, 18, 86, 23), and the first and last columns two middle numbers (88, 10, 25, 16) all add up to the sum of 139.
Multiplicative magic squares
Instead of adding the numbers in each row, column and diagonal, one can apply some other operation. For example, a multiplicative magic square has a constant product of numbers. A multiplicative magic square can be derived from an additive magic square by raising 2 (or any other integer) to the power of each element, because the logarithm of the product of 2 numbers is the sum of logarithm of each. Alternatively, if any 3 numbers in a line are 2a, 2b and 2c, their product is 2a+b+c, which is constant if a+b+c is constant, as they would be if a, b and c were taken from ordinary (additive) magic square. For example, the original Lo-Shu magic square becomes:
Other examples of multiplicative magic squares include:
Multiplicative magic squares of complex numbers
Still using Ali Skalli's non iterative method, it is possible to produce an infinity of multiplicative magic squares of complex numbers belonging to set. On the example below, the real and imaginary parts are integer numbers, but they can also belong to the entire set of real numbers .
The product is: −352,507,340,640 − 400,599,719,520 i.
Additive-multiplicative magic and semimagic squares
Additive-multiplicative magic squares and semimagic squares satisfy properties of both ordinary and multiplicative magic squares and semimagic squares, respectively.
It is unknown if any additive-multiplicative magic squares smaller than 7×7 exist, but it has been proven that no 3×3 or 4×4 additive-multiplicative magic squares and no 3×3 additive-multiplicative semimagic squares exist.
Geometric magic squares
Magic squares may be constructed which contain geometric shapes instead of numbers. Such squares, known as geometric magic squares, were invented and named by Lee Sallows in 2001.
In the example shown the shapes appearing are two dimensional. It was Sallows' discovery that all magic squares are geometric, the numbers that appear in numerical magic squares can be interpreted as a shorthand notation which indicates the lengths of straight line segments that are the geometric 'shapes' occurring in the square. That is, numerical magic squares are that special case of a geometric magic square using one dimensional shapes.
Area magic squares
In 2017, following initial ideas of William Walkington and Inder Taneja, the first linear area magic square (L-AMS) was constructed by Walter Trump.
Other magic shapes
Other two dimensional shapes than squares can be considered. The general case is to consider a design with N parts to be magic if the N parts are labeled with the numbers 1 through N and a number of identical sub-designs give the same sum. Examples include magic circles, magic rectangles, magic triangles magic stars, magic hexagons, magic diamonds. Going up in dimension results in magic spheres, magic cylinders, magic cubes, magic parallelepiped, magic solids, and other magic hypercubes.
Possible magic shapes are constrained by the number of equal-sized, equal-sum subsets of the chosen set of labels. For example, if one proposes to form a magic shape labeling the parts with {1, 2, 3, 4}, the sub-designs will have to be labeled with {1,4} and {2,3}.
Related problems
n-Queens problem
In 1992, Demirörs, Rafraf, and Tanik published a method for converting some magic squares into n-queens solutions, and vice versa.
Magic squares in occultism
Magic squares of order 3 through 9, assigned to the seven planets, and described as means to attract the influence of planets and their angels (or demons) during magical practices, can be found in several manuscripts all around Europe starting at least since the 15th century. Among the best known, the Liber de Angelis, a magical handbook written around 1440, is included in Cambridge Univ. Lib. MS Dd.xi.45. The text of the Liber de Angelis is very close to that of De septem quadraturis planetarum seu quadrati magici, another handbook of planetary image magic contained in the Codex 793 of the Biblioteka Jagiellońska (Ms BJ 793). The magical operations involve engraving the appropriate square on a plate made with the metal assigned to the corresponding planet, as well as performing a variety of rituals. For instance, the 3×3 square, that belongs to Saturn, has to be inscribed on a lead plate. It will, in particular, help women during a difficult childbirth.
In about 1510 Heinrich Cornelius Agrippa wrote De Occulta Philosophia, drawing on the Hermetic and magical works of Marsilio Ficino and Pico della Mirandola. In its 1531 edition, he expounded on the magical virtues of the seven magical squares of orders 3 to 9, each associated with one of the astrological planets, much in the same way as the older texts did. This book was very influential throughout Europe until the Counter-Reformation, and Agrippa's magic squares, sometimes called kameas, continue to be used within modern ceremonial magic in much the same way as he first prescribed.
The most common use for these kameas is to provide a pattern upon which to construct the sigils of spirits, angels or demons; the letters of the entity's name are converted into numbers, and lines are traced through the pattern that these successive numbers make on the kamea.
In a magical context, the term magic square is also applied to a variety of word squares or number squares found in magical grimoires, including some that do not follow any obvious pattern, and even those with differing numbers of rows and columns. They are generally intended for use as talismans. For instance the following squares are: The Sator square, one of the most famous magic squares found in a number of grimoires including the Key of Solomon; a square "to overcome envy", from The Book of Power; and two squares from The Book of the Sacred Magic of Abramelin the Mage, the first to cause the illusion of a superb palace to appear, and the second to be worn on the head of a child during an angelic invocation:
Magic squares in popular culture
In Goethe's Faust, the witch's spell used to make a youth elixir for Faust, the , has been interpreted as a construction of a magic square.
The English composer Peter Maxwell Davies has used magic squares to structure many of his compositions. For example, his 1975 Ave Maris Stella uses the 9×9 magic square of Moon while his 1977 A Mirror of Whitening Light uses the 8×8 magic square of Mercury to create the entire set of notes and durations for the piece. His other works that employ magic squares include The Lighthouse (1979), Resurrection (1987), Strathclyde Concerto No. 3 for Horn and Trumpet (1989), as well as many of his symphonies. According to Davies' own account:
A magic square in a musical composition is not a block of numbers – it is a generating principle, to be learned and known intimately, perceived inwardly as a multi-dimensional projection into that vast (chaotic!) area of the internal ear – the space/time crucible – where music is conceived. ... Projected onto the page, a magic square is a dead, black conglomeration of digits; tune in, and one hears a powerful, orbiting dynamo of musical images, glowing with numen and lumen.
Magic squares, including Benjamin Franklin's, appear as clues to the mystery in Katherine Neville's novels The Eight and The Fire.
Magic squares play a role in Steve Martin's 2003 novel The Pleasure of My Company.
Dürer's magic square and his Melencolia I both also played large roles in Dan Brown's 2009 novel, The Lost Symbol.
In the 2011 Korean television drama Deep Rooted Tree, King Sejong is shown attempting to construct a 33×33 magic square using lunch boxes. He ultimately discovers the "pyramid method" and completes the magic square with the help of an army of court attendants. This inspires him to create a more just form of government ruled by reason and words rather than military might.
On October 9, 2014, the post office of Macao in the People's Republic of China issued a series of stamps based on magic squares. The figure below shows the stamps featuring the nine magic squares chosen to be in this collection.
The metallic artifact at the center of The X-Files episode "Biogenesis" is alleged by Chuck Burks to be a magic square.
Mathematician Matt Parker attempted to create a 3×3 magic square using square numbers in a YouTube video on the Numberphile channel. His failed attempt is known as the Parker square.
The first season Stargate Atlantis episode "Brotherhood" involves completing a magic square as part of a puzzle guarding a powerful Ancient artefact.
Magic Squares are also featured in the 2019 Spanish film Vivir dos veces.
| Mathematics | Sums and products | null |
52247 | https://en.wikipedia.org/wiki/Fourier%20transform | Fourier transform | In mathematics, the Fourier transform (FT) is an integral transform that takes a function as input and outputs another function that describes the extent to which various frequencies are present in the original function. The output of the transform is a complex-valued function of frequency. The term Fourier transform refers to both this complex-valued function and the mathematical operation. When a distinction needs to be made, the output of the operation is sometimes called the frequency domain representation of the original function. The Fourier transform is analogous to decomposing the sound of a musical chord into the intensities of its constituent pitches.
Functions that are localized in the time domain have Fourier transforms that are spread out across the frequency domain and vice versa, a phenomenon known as the uncertainty principle. The critical case for this principle is the Gaussian function, of substantial importance in probability theory and statistics as well as in the study of physical phenomena exhibiting normal distribution (e.g., diffusion). The Fourier transform of a Gaussian function is another Gaussian function. Joseph Fourier introduced sine and cosine transforms (which correspond to the imaginary and real components of the modern Fourier transform) in his study of heat transfer, where Gaussian functions appear as solutions of the heat equation.
The Fourier transform can be formally defined as an improper Riemann integral, making it an integral transform, although this definition is not suitable for many applications requiring a more sophisticated integration theory. For example, many relatively simple applications use the Dirac delta function, which can be treated formally as if it were a function, but the justification requires a mathematically more sophisticated viewpoint.
The Fourier transform can also be generalized to functions of several variables on Euclidean space, sending a function of 'position space' to a function of momentum (or a function of space and time to a function of 4-momentum). This idea makes the spatial Fourier transform very natural in the study of waves, as well as in quantum mechanics, where it is important to be able to represent wave solutions as functions of either position or momentum and sometimes both. In general, functions to which Fourier methods are applicable are complex-valued, and possibly vector-valued. Still further generalization is possible to functions on groups, which, besides the original Fourier transform on or , notably includes the discrete-time Fourier transform (DTFT, group = ), the discrete Fourier transform (DFT, group = ) and the Fourier series or circular Fourier transform (group = , the unit circle ≈ closed finite interval with endpoints identified). The latter is routinely employed to handle periodic functions. The fast Fourier transform (FFT) is an algorithm for computing the DFT.
Definition
The Fourier transform of a complex-valued (Lebesgue) integrable function on the real line, is the complex valued function , defined by the integral
Evaluating the Fourier transform for all values of produces the frequency-domain function, and it converges at all frequencies to a continuous function tending to zero at infinity. If decays with all derivatives, i.e.,
then converges for all frequencies and, by the Riemann–Lebesgue lemma, also decays with all derivatives.
First introduced in Fourier's Analytical Theory of Heat., the corresponding inversion formula for "sufficiently nice" functions is given by the Fourier inversion theorem, i.e.,
The functions and are referred to as a Fourier transform pair. A common notation for designating transform pairs is:
for example
By analogy, the Fourier series can be regarded as abstract Fourier transform on the group of integers. That is, the synthesis of a sequence of complex numbers is defined by the Fourier transform
such that are given by the inversion formula, i.e., the analysis
for some complex-valued, -periodic function defined on a bounded interval . When the constituent frequencies are a continuum: and .
In other words, on the finite interval the function has a discrete decomposition in the periodic functions . On the infinite interval the function has a continuous decomposition in periodic functions .
Lebesgue integrable functions
A measurable function is called (Lebesgue) integrable if the Lebesgue integral of its absolute value is finite:
If is Lebesgue integrable then the Fourier transform, given by , is well-defined for all . Furthermore, is bounded, uniformly continuous and (by the Riemann–Lebesgue lemma) zero at infinity.
The space is the space of measurable functions for which the norm is finite, modulo the equivalence relation of equality almost everywhere. The Fourier transform is one-to-one on . However, there is no easy characterization of the image, and thus no easy characterization of the inverse transform. In particular, is no longer valid, as it was stated only under the hypothesis that decayed with all derivatives.
While defines the Fourier transform for (complex-valued) functions in , it is not well-defined for other integrability classes, most importantly the space of square-integrable functions . For example, the function is in but not and therefore the Lebesgue integral does not exist. However, the Fourier transform on the dense subspace admits a unique continuous extension to a unitary operator on . This extension is important in part because, unlike the case of , the Fourier transform is an automorphism of the space .
In such cases, the Fourier transform can be obtained explicitly by regularizing the integral, and then passing to a limit. In practice, the integral is often regarded as an improper integral instead of a proper Lebesgue integral, but sometimes for convergence one needs to use weak limit or principal value instead of the (pointwise) limits implicit in an improper integral. and each gives three rigorous ways of extending the Fourier transform to square integrable functions using this procedure. A general principle in working with the Fourier transform is that Gaussians are dense in , and the various features of the Fourier transform, such as its unitarity, are easily inferred for Gaussians. Many of the properties of the Fourier transform, can then be proven from two facts about Gaussians:
that is its own Fourier transform; and
that the Gaussian integral
A feature of the Fourier transform is that it is a homomorphism of Banach algebras from equipped with the convolution operation to the Banach algebra of continuous functions under the (supremum) norm. The conventions chosen in this article are those of harmonic analysis, and are characterized as the unique conventions such that the Fourier transform is both unitary on and an algebra homomorphism from to , without renormalizing the Lebesgue measure.
Angular frequency (ω)
When the independent variable () represents time (often denoted by ), the transform variable () represents frequency (often denoted by ). For example, if time is measured in seconds, then frequency is in hertz. The Fourier transform can also be written in terms of angular frequency, whose units are radians per second.
The substitution into produces this convention, where function is relabeled
Unlike the definition, the Fourier transform is no longer a unitary transformation, and there is less symmetry between the formulas for the transform and its inverse. Those properties are restored by splitting the factor evenly between the transform and its inverse, which leads to another convention:
Variations of all three conventions can be created by conjugating the complex-exponential kernel of both the forward and the reverse transform. The signs must be opposites.
Background
History
In 1822, Fourier claimed (see ) that any function, whether continuous or discontinuous, can be expanded into a series of sines. That important work was corrected and expanded upon by others to provide the foundation for the various forms of the Fourier transform used since.
Complex sinusoids
In general, the coefficients are complex numbers, which have two equivalent forms (see Euler's formula):
The product with () has these forms:
which conveys both amplitude and phase of frequency Likewise, the intuitive interpretation of is that multiplying by has the effect of subtracting from every frequency component of function Only the component that was at frequency can produce a non-zero value of the infinite integral, because (at least formally) all the other shifted components are oscillatory and integrate to zero. (see )
It is noteworthy how easily the product was simplified using the polar form, and how easily the rectangular form was deduced by an application of Euler's formula.
Negative frequency
Euler's formula introduces the possibility of negative And is defined Only certain complex-valued have transforms (See Analytic signal. A simple example is ) But negative frequency is necessary to characterize all other complex-valued found in signal processing, partial differential equations, radar, nonlinear optics, quantum mechanics, and others.
For a real-valued has the symmetry property (see below). This redundancy enables to distinguish from But of course it cannot tell us the actual sign of because and are indistinguishable on just the real numbers line.
Fourier transform for periodic functions
The Fourier transform of a periodic function cannot be defined using the integral formula directly. In order for integral in to be defined the function must be absolutely integrable. Instead it is common to use Fourier series. It is possible to extend the definition to include periodic functions by viewing them as tempered distributions.
This makes it possible to see a connection between the Fourier series and the Fourier transform for periodic functions that have a convergent Fourier series. If is a periodic function, with period , that has a convergent Fourier series, then:
where are the Fourier series coefficients of , and is the Dirac delta function. In other words, the Fourier transform is a Dirac comb function whose teeth are multiplied by the Fourier series coefficients.
Sampling the Fourier transform
The Fourier transform of an integrable function can be sampled at regular intervals of arbitrary length These samples can be deduced from one cycle of a periodic function which has Fourier series coefficients proportional to those samples by the Poisson summation formula:
The integrability of ensures the periodic summation converges. Therefore, the samples can be determined by Fourier series analysis:
When has compact support, has a finite number of terms within the interval of integration. When does not have compact support, numerical evaluation of requires an approximation, such as tapering or truncating the number of terms.
Units
The frequency variable must have inverse units to the units of the original function's domain (typically named or ). For example, if is measured in seconds, should be in cycles per second or hertz. If the scale of time is in units of seconds, then another Greek letter is typically used instead to represent angular frequency (where ) in units of radians per second. If using for units of length, then must be in inverse length, e.g., wavenumbers. That is to say, there are two versions of the real line: one which is the range of and measured in units of and the other which is the range of and measured in inverse units to the units of These two distinct versions of the real line cannot be equated with each other. Therefore, the Fourier transform goes from one space of functions to a different space of functions: functions which have a different domain of definition.
In general, must always be taken to be a linear form on the space of its domain, which is to say that the second real line is the dual space of the first real line. See the article on linear algebra for a more formal explanation and for more details. This point of view becomes essential in generalizations of the Fourier transform to general symmetry groups, including the case of Fourier series.
That there is no one preferred way (often, one says "no canonical way") to compare the two versions of the real line which are involved in the Fourier transform—fixing the units on one line does not force the scale of the units on the other line—is the reason for the plethora of rival conventions on the definition of the Fourier transform. The various definitions resulting from different choices of units differ by various constants.
In other conventions, the Fourier transform has in the exponent instead of , and vice versa for the inversion formula. This convention is common in modern physics and is the default for Wolfram Alpha, and does not mean that the frequency has become negative, since there is no canonical definition of positivity for frequency of a complex wave. It simply means that is the amplitude of the wave instead of the wave (the former, with its minus sign, is often seen in the time dependence for Sinusoidal plane-wave solutions of the electromagnetic wave equation, or in the time dependence for quantum wave functions). Many of the identities involving the Fourier transform remain valid in those conventions, provided all terms that explicitly involve have it replaced by . In Electrical engineering the letter is typically used for the imaginary unit instead of because is used for current.
When using dimensionless units, the constant factors might not even be written in the transform definition. For instance, in probability theory, the characteristic function of the probability density function of a random variable of continuous type is defined without a negative sign in the exponential, and since the units of are ignored, there is no 2 either:
(In probability theory, and in mathematical statistics, the use of the Fourier—Stieltjes transform is preferred, because so many random variables are not of continuous type, and do not possess a density function, and one must treat not functions but distributions, i.e., measures which possess "atoms".)
From the higher point of view of group characters, which is much more abstract, all these arbitrary choices disappear, as will be explained in the later section of this article, which treats the notion of the Fourier transform of a function on a locally compact Abelian group.
Properties
Let and represent integrable functions Lebesgue-measurable on the real line satisfying:
We denote the Fourier transforms of these functions as and respectively.
Basic properties
The Fourier transform has the following basic properties:
Linearity
Time shifting
Frequency shifting
Time scaling
The case leads to the time-reversal property:
Symmetry
When the real and imaginary parts of a complex function are decomposed into their even and odd parts, there are four components, denoted below by the subscripts RE, RO, IE, and IO. And there is a one-to-one mapping between the four components of a complex time function and the four components of its complex frequency transform:
From this, various relationships are apparent, for example:
The transform of a real-valued function is the conjugate symmetric function Conversely, a conjugate symmetric transform implies a real-valued time-domain.
The transform of an imaginary-valued function is the conjugate antisymmetric function and the converse is true.
The transform of a conjugate symmetric function is the real-valued function and the converse is true.
The transform of a conjugate antisymmetric function is the imaginary-valued function and the converse is true.
Conjugation
(Note: the ∗ denotes complex conjugation.)
In particular, if is real, then is even symmetric (aka Hermitian function):
And if is purely imaginary, then is odd symmetric:
Real and imaginary parts
Zero frequency component
Substituting in the definition, we obtain:
The integral of over its domain is known as the average value or DC bias of the function.
Uniform continuity and the Riemann–Lebesgue lemma
The Fourier transform may be defined in some cases for non-integrable functions, but the Fourier transforms of integrable functions have several strong properties.
The Fourier transform of any integrable function is uniformly continuous and
By the Riemann–Lebesgue lemma,
However, need not be integrable. For example, the Fourier transform of the rectangular function, which is integrable, is the sinc function, which is not Lebesgue integrable, because its improper integrals behave analogously to the alternating harmonic series, in converging to a sum without being absolutely convergent.
It is not generally possible to write the inverse transform as a Lebesgue integral. However, when both and are integrable, the inverse equality
holds for almost every . As a result, the Fourier transform is injective on .
Plancherel theorem and Parseval's theorem
Let and be integrable, and let and be their Fourier transforms. If and are also square-integrable, then the Parseval formula follows:
where the bar denotes complex conjugation.
The Plancherel theorem, which follows from the above, states that
Plancherel's theorem makes it possible to extend the Fourier transform, by a continuity argument, to a unitary operator on . On , this extension agrees with original Fourier transform defined on , thus enlarging the domain of the Fourier transform to (and consequently to for ). Plancherel's theorem has the interpretation in the sciences that the Fourier transform preserves the energy of the original quantity. The terminology of these formulas is not quite standardised. Parseval's theorem was proved only for Fourier series, and was first proved by Lyapunov. But Parseval's formula makes sense for the Fourier transform as well, and so even though in the context of the Fourier transform it was proved by Plancherel, it is still often referred to as Parseval's formula, or Parseval's relation, or even Parseval's theorem.
See Pontryagin duality for a general formulation of this concept in the context of locally compact abelian groups.
Convolution theorem
The Fourier transform translates between convolution and multiplication of functions. If and are integrable functions with Fourier transforms and respectively, then the Fourier transform of the convolution is given by the product of the Fourier transforms and (under other conventions for the definition of the Fourier transform a constant factor may appear).
This means that if:
where denotes the convolution operation, then:
In linear time invariant (LTI) system theory, it is common to interpret as the impulse response of an LTI system with input and output , since substituting the unit impulse for yields . In this case, represents the frequency response of the system.
Conversely, if can be decomposed as the product of two square integrable functions and , then the Fourier transform of is given by the convolution of the respective Fourier transforms and .
Cross-correlation theorem
In an analogous manner, it can be shown that if is the cross-correlation of and :
then the Fourier transform of is:
As a special case, the autocorrelation of function is:
for which
Differentiation
Suppose is an absolutely continuous differentiable function, and both and its derivative are integrable. Then the Fourier transform of the derivative is given by
More generally, the Fourier transformation of the th derivative is given by
Analogously, , so
By applying the Fourier transform and using these formulas, some ordinary differential equations can be transformed into algebraic equations, which are much easier to solve. These formulas also give rise to the rule of thumb " is smooth if and only if quickly falls to 0 for ." By using the analogous rules for the inverse Fourier transform, one can also say " quickly falls to 0 for if and only if is smooth."
Eigenfunctions
The Fourier transform is a linear transform which has eigenfunctions obeying with
A set of eigenfunctions is found by noting that the homogeneous differential equation
leads to eigenfunctions of the Fourier transform as long as the form of the equation remains invariant under Fourier transform. In other words, every solution and its Fourier transform obey the same equation. Assuming uniqueness of the solutions, every solution must therefore be an eigenfunction of the Fourier transform. The form of the equation remains unchanged under Fourier transform if can be expanded in a power series in which for all terms the same factor of either one of arises from the factors introduced by the differentiation rules upon Fourier transforming the homogeneous differential equation because this factor may then be cancelled. The simplest allowable leads to the standard normal distribution.
More generally, a set of eigenfunctions is also found by noting that the differentiation rules imply that the ordinary differential equation
with constant and being a non-constant even function remains invariant in form when applying the Fourier transform to both sides of the equation. The simplest example is provided by which is equivalent to considering the Schrödinger equation for the quantum harmonic oscillator. The corresponding solutions provide an important choice of an orthonormal basis for and are given by the "physicist's" Hermite functions. Equivalently one may use
where are the "probabilist's" Hermite polynomials, defined as
Under this convention for the Fourier transform, we have that
In other words, the Hermite functions form a complete orthonormal system of eigenfunctions for the Fourier transform on . However, this choice of eigenfunctions is not unique. Because of there are only four different eigenvalues of the Fourier transform (the fourth roots of unity ±1 and ±) and any linear combination of eigenfunctions with the same eigenvalue gives another eigenfunction. As a consequence of this, it is possible to decompose as a direct sum of four spaces , , , and where the Fourier transform acts on simply by multiplication by .
Since the complete set of Hermite functions provides a resolution of the identity they diagonalize the Fourier operator, i.e. the Fourier transform can be represented by such a sum of terms weighted by the above eigenvalues, and these sums can be explicitly summed:
This approach to define the Fourier transform was first proposed by Norbert Wiener. Among other properties, Hermite functions decrease exponentially fast in both frequency and time domains, and they are thus used to define a generalization of the Fourier transform, namely the fractional Fourier transform used in time–frequency analysis. In physics, this transform was introduced by Edward Condon. This change of basis functions becomes possible because the Fourier transform is a unitary transform when using the right conventions. Consequently, under the proper conditions it may be expected to result from a self-adjoint generator via
The operator is the number operator of the quantum harmonic oscillator written as
It can be interpreted as the generator of fractional Fourier transforms for arbitrary values of , and of the conventional continuous Fourier transform for the particular value with the Mehler kernel implementing the corresponding active transform. The eigenfunctions of are the Hermite functions which are therefore also eigenfunctions of
Upon extending the Fourier transform to distributions the Dirac comb is also an eigenfunction of the Fourier transform.
Inversion and periodicity
Under suitable conditions on the function , it can be recovered from its Fourier transform . Indeed, denoting the Fourier transform operator by , so , then for suitable functions, applying the Fourier transform twice simply flips the function: , which can be interpreted as "reversing time". Since reversing time is two-periodic, applying this twice yields , so the Fourier transform operator is four-periodic, and similarly the inverse Fourier transform can be obtained by applying the Fourier transform three times: . In particular the Fourier transform is invertible (under suitable conditions).
More precisely, defining the parity operator such that , we have:
These equalities of operators require careful definition of the space of functions in question, defining equality of functions (equality at every point? equality almost everywhere?) and defining equality of operators – that is, defining the topology on the function space and operator space in question. These are not true for all functions, but are true under various conditions, which are the content of the various forms of the Fourier inversion theorem.
This fourfold periodicity of the Fourier transform is similar to a rotation of the plane by 90°, particularly as the two-fold iteration yields a reversal, and in fact this analogy can be made precise. While the Fourier transform can simply be interpreted as switching the time domain and the frequency domain, with the inverse Fourier transform switching them back, more geometrically it can be interpreted as a rotation by 90° in the time–frequency domain (considering time as the -axis and frequency as the -axis), and the Fourier transform can be generalized to the fractional Fourier transform, which involves rotations by other angles. This can be further generalized to linear canonical transformations, which can be visualized as the action of the special linear group on the time–frequency plane, with the preserved symplectic form corresponding to the uncertainty principle, below. This approach is particularly studied in signal processing, under time–frequency analysis.
Connection with the Heisenberg group
The Heisenberg group is a certain group of unitary operators on the Hilbert space of square integrable complex valued functions on the real line, generated by the translations and multiplication by , . These operators do not commute, as their (group) commutator is
which is multiplication by the constant (independent of ) (the circle group of unit modulus complex numbers). As an abstract group, the Heisenberg group is the three-dimensional Lie group of triples , with the group law
Denote the Heisenberg group by . The above procedure describes not only the group structure, but also a standard unitary representation of on a Hilbert space, which we denote by . Define the linear automorphism of by
so that . This can be extended to a unique automorphism of :
According to the Stone–von Neumann theorem, the unitary representations and are unitarily equivalent, so there is a unique intertwiner such that
This operator is the Fourier transform.
Many of the standard properties of the Fourier transform are immediate consequences of this more general framework. For example, the square of the Fourier transform, , is an intertwiner associated with , and so we have is the reflection of the original function .
Complex domain
The integral for the Fourier transform
can be studied for complex values of its argument . Depending on the properties of , this might not converge off the real axis at all, or it might converge to a complex analytic function for all values of , or something in between.
The Paley–Wiener theorem says that is smooth (i.e., -times differentiable for all positive integers ) and compactly supported if and only if is a holomorphic function for which there exists a constant such that for any integer ,
for some constant . (In this case, is supported on .) This can be expressed by saying that is an entire function which is rapidly decreasing in (for fixed ) and of exponential growth in (uniformly in ).
(If is not smooth, but only , the statement still holds provided .) The space of such functions of a complex variable is called the Paley—Wiener space. This theorem has been generalised to semisimple Lie groups.
If is supported on the half-line , then is said to be "causal" because the impulse response function of a physically realisable filter must have this property, as no effect can precede its cause. Paley and Wiener showed that then extends to a holomorphic function on the complex lower half-plane which tends to zero as goes to infinity. The converse is false and it is not known how to characterise the Fourier transform of a causal function.
Laplace transform
The Fourier transform is related to the Laplace transform , which is also used for the solution of differential equations and the analysis of filters.
It may happen that a function for which the Fourier integral does not converge on the real axis at all, nevertheless has a complex Fourier transform defined in some region of the complex plane.
For example, if is of exponential growth, i.e.,
for some constants , then
convergent for all , is the two-sided Laplace transform of .
The more usual version ("one-sided") of the Laplace transform is
If is also causal, and analytical, then: Thus, extending the Fourier transform to the complex domain means it includes the Laplace transform as a special case in the case of causal functions—but with the change of variable .
From another, perhaps more classical viewpoint, the Laplace transform by its form involves an additional exponential regulating term which lets it converge outside of the imaginary line where the Fourier transform is defined. As such it can converge for at most exponentially divergent series and integrals, whereas the original Fourier decomposition cannot, enabling analysis of systems with divergent or critical elements. Two particular examples from linear signal processing are the construction of allpass filter networks from critical comb and mitigating filters via exact pole-zero cancellation on the unit circle. Such designs are common in audio processing, where highly nonlinear phase response is sought for, as in reverb.
Furthermore, when extended pulselike impulse responses are sought for signal processing work, the easiest way to produce them is to have one circuit which produces a divergent time response, and then to cancel its divergence through a delayed opposite and compensatory response. There, only the delay circuit in-between admits a classical Fourier description, which is critical. Both the circuits to the side are unstable, and do not admit a convergent Fourier decomposition. However, they do admit a Laplace domain description, with identical half-planes of convergence in the complex plane (or in the discrete case, the Z-plane), wherein their effects cancel.
In modern mathematics the Laplace transform is conventionally subsumed under the aegis Fourier methods. Both of them are subsumed by the far more general, and more abstract, idea of harmonic analysis.
Inversion
Still with , if is complex analytic for , then
by Cauchy's integral theorem. Therefore, the Fourier inversion formula can use integration along different lines, parallel to the real axis.
Theorem: If for , and for some constants , then
for any .
This theorem implies the Mellin inversion formula for the Laplace transformation,
for any , where is the Laplace transform of .
The hypotheses can be weakened, as in the results of Carleson and Hunt, to being , provided that be of bounded variation in a closed neighborhood of (cf. Dini test), the value of at be taken to be the arithmetic mean of the left and right limits, and that the integrals be taken in the sense of Cauchy principal values.
versions of these inversion formulas are also available.
Fourier transform on Euclidean space
The Fourier transform can be defined in any arbitrary number of dimensions . As with the one-dimensional case, there are many conventions. For an integrable function , this article takes the definition:
where and are -dimensional vectors, and is the dot product of the vectors. Alternatively, can be viewed as belonging to the dual vector space , in which case the dot product becomes the contraction of and , usually written as .
All of the basic properties listed above hold for the -dimensional Fourier transform, as do Plancherel's and Parseval's theorem. When the function is integrable, the Fourier transform is still uniformly continuous and the Riemann–Lebesgue lemma holds.
Uncertainty principle
Generally speaking, the more concentrated is, the more spread out its Fourier transform must be. In particular, the scaling property of the Fourier transform may be seen as saying: if we squeeze a function in , its Fourier transform stretches out in . It is not possible to arbitrarily concentrate both a function and its Fourier transform.
The trade-off between the compaction of a function and its Fourier transform can be formalized in the form of an uncertainty principle by viewing a function and its Fourier transform as conjugate variables with respect to the symplectic form on the time–frequency domain: from the point of view of the linear canonical transformation, the Fourier transform is rotation by 90° in the time–frequency domain, and preserves the symplectic form.
Suppose is an integrable and square-integrable function. Without loss of generality, assume that is normalized:
It follows from the Plancherel theorem that is also normalized.
The spread around may be measured by the dispersion about zero defined by
In probability terms, this is the second moment of about zero.
The uncertainty principle states that, if is absolutely continuous and the functions and are square integrable, then
The equality is attained only in the case
where is arbitrary and so that is -normalized. In other words, where is a (normalized) Gaussian function with variance , centered at zero, and its Fourier transform is a Gaussian function with variance .
In fact, this inequality implies that:
for any , .
In quantum mechanics, the momentum and position wave functions are Fourier transform pairs, up to a factor of the Planck constant. With this constant properly taken into account, the inequality above becomes the statement of the Heisenberg uncertainty principle.
A stronger uncertainty principle is the Hirschman uncertainty principle, which is expressed as:
where is the differential entropy of the probability density function :
where the logarithms may be in any base that is consistent. The equality is attained for a Gaussian, as in the previous case.
Sine and cosine transforms
Fourier's original formulation of the transform did not use complex numbers, but rather sines and cosines. Statisticians and others still use this form. An absolutely integrable function for which Fourier inversion holds can be expanded in terms of genuine frequencies (avoiding negative frequencies, which are sometimes considered hard to interpret physically) by
This is called an expansion as a trigonometric integral, or a Fourier integral expansion. The coefficient functions and can be found by using variants of the Fourier cosine transform and the Fourier sine transform (the normalisations are, again, not standardised):
and
Older literature refers to the two transform functions, the Fourier cosine transform, , and the Fourier sine transform, .
The function can be recovered from the sine and cosine transform using
together with trigonometric identities. This is referred to as Fourier's integral formula.
Spherical harmonics
Let the set of homogeneous harmonic polynomials of degree on be denoted by . The set consists of the solid spherical harmonics of degree . The solid spherical harmonics play a similar role in higher dimensions to the Hermite polynomials in dimension one. Specifically, if for some in , then . Let the set be the closure in of linear combinations of functions of the form where is in . The space is then a direct sum of the spaces and the Fourier transform maps each space to itself and is possible to characterize the action of the Fourier transform on each space .
Let (with in ), then
where
Here denotes the Bessel function of the first kind with order . When this gives a useful formula for the Fourier transform of a radial function. This is essentially the Hankel transform. Moreover, there is a simple recursion relating the cases and allowing to compute, e.g., the three-dimensional Fourier transform of a radial function from the one-dimensional one.
Restriction problems
In higher dimensions it becomes interesting to study restriction problems for the Fourier transform. The Fourier transform of an integrable function is continuous and the restriction of this function to any set is defined. But for a square-integrable function the Fourier transform could be a general class of square integrable functions. As such, the restriction of the Fourier transform of an function cannot be defined on sets of measure 0. It is still an active area of study to understand restriction problems in for . It is possible in some cases to define the restriction of a Fourier transform to a set , provided has non-zero curvature. The case when is the unit sphere in is of particular interest. In this case the Tomas–Stein restriction theorem states that the restriction of the Fourier transform to the unit sphere in is a bounded operator on provided .
One notable difference between the Fourier transform in 1 dimension versus higher dimensions concerns the partial sum operator. Consider an increasing collection of measurable sets indexed by : such as balls of radius centered at the origin, or cubes of side . For a given integrable function , consider the function defined by:
Suppose in addition that . For and , if one takes , then converges to in as tends to infinity, by the boundedness of the Hilbert transform. Naively one may hope the same holds true for . In the case that is taken to be a cube with side length , then convergence still holds. Another natural candidate is the Euclidean ball . In order for this partial sum operator to converge, it is necessary that the multiplier for the unit ball be bounded in . For it is a celebrated theorem of Charles Fefferman that the multiplier for the unit ball is never bounded unless . In fact, when , this shows that not only may fail to converge to in , but for some functions , is not even an element of .
Fourier transform on function spaces
The definition of the Fourier transform naturally extends from to . That is, if then the Fourier transform
is given by
This operator is bounded as
which shows that its operator norm is bounded by . The Riemann–Lebesgue lemma shows that if then its Fourier transform actually belongs to the space of continuous functions which vanish at infinity, i.e., . Furthermore, the image of under is a strict subset of .
Similarly to the case of one variable, the Fourier transform can be defined on . The Fourier transform in is no longer given by an ordinary Lebesgue integral, although it can be computed by an improper integral, i.e.,
where the limit is taken in the sense.
Furthermore, is a unitary operator. For an operator to be unitary it is sufficient to show that it is bijective and preserves the inner product, so in this case these follow from the Fourier inversion theorem combined with the fact that for any we have
In particular, the image of is itself under the Fourier transform.
On other Lp
For , the Fourier transform can be defined on by Marcinkiewicz interpolation, which amounts to decomposing such functions into a fat tail part in plus a fat body part in . In each of these spaces, the Fourier transform of a function in is in , where is the Hölder conjugate of (by the Hausdorff–Young inequality). However, except for , the image is not easily characterized. Further extensions become more technical. The Fourier transform of functions in for the range requires the study of distributions. In fact, it can be shown that there are functions in with so that the Fourier transform is not defined as a function.
Tempered distributions
One might consider enlarging the domain of the Fourier transform from by considering generalized functions, or distributions. A distribution on is a continuous linear functional on the space of compactly supported smooth functions (i.e. bump functions), equipped with a suitable topology. Since is dense in , the Plancherel theorem allows one to extend the definition of the Fourier transform to general functions in by continuity arguments. The strategy is then to consider the action of the Fourier transform on and pass to distributions by duality. The obstruction to doing this is that the Fourier transform does not map to . In fact the Fourier transform of an element in can not vanish on an open set; see the above discussion on the uncertainty principle.
The Fourier transform can also be defined for tempered distributions , dual to the space of Schwartz functions . A Schwartz function is a smooth function that decays at infinity, along with all of its derivatives, hence and:
The Fourier transform is an automorphism of the Schwartz space and, by duality, also an automorphism of the space of tempered distributions. The tempered distributions include well-behaved functions of polynomial growth, distributions of compact support as well as all the integrable functions mentioned above.
For the definition of the Fourier transform of a tempered distribution, let and be integrable functions, and let and be their Fourier transforms respectively. Then the Fourier transform obeys the following multiplication formula,
Every integrable function defines (induces) a distribution by the relation
So it makes sense to define the Fourier transform of a tempered distribution by the duality:
Extending this to all tempered distributions gives the general definition of the Fourier transform.
Distributions can be differentiated and the above-mentioned compatibility of the Fourier transform with differentiation and convolution remains true for tempered distributions.
Generalizations
Fourier–Stieltjes transform on measurable spaces
The Fourier transform of a finite Borel measure on is given by the continuous function:
and called the Fourier-Stieltjes transform due to its connection with the Riemann-Stieltjes integral representation of (Radon) measures. If is the probability distribution of a random variable then its Fourier–Stieltjes transform is, by definition, a characteristic function. If, in addition, the probability distribution has a probability density function, this definition is subject to the usual Fourier transform. Stated more generally, when is absolutely continuous with respect to the Lebesgue measure, i.e.,
then
and the Fourier-Stieltjes transform reduces to the usual definition of the Fourier transform. That is, the notable difference with the Fourier transform of integrable functions is that the Fourier-Stieltjes transform need not vanish at infinity, i.e., the Riemann–Lebesgue lemma fails for measures.
Bochner's theorem characterizes which functions may arise as the Fourier–Stieltjes transform of a positive measure on the circle.
One example of a finite Borel measure that is not a function is the Dirac measure. Its Fourier transform is a constant function (whose value depends on the form of the Fourier transform used).
Locally compact abelian groups
The Fourier transform may be generalized to any locally compact abelian group, i.e., an abelian group that is also a locally compact Hausdorff space such that the group operation is continuous. If is a locally compact abelian group, it has a translation invariant measure , called Haar measure. For a locally compact abelian group , the set of irreducible, i.e. one-dimensional, unitary representations are called its characters. With its natural group structure and the topology of uniform convergence on compact sets (that is, the topology induced by the compact-open topology on the space of all continuous functions from to the circle group), the set of characters is itself a locally compact abelian group, called the Pontryagin dual of . For a function in , its Fourier transform is defined by
The Riemann–Lebesgue lemma holds in this case; is a function vanishing at infinity on .
The Fourier transform on is an example; here is a locally compact abelian group, and the Haar measure on can be thought of as the Lebesgue measure on [0,1). Consider the representation of on the complex plane that is a 1-dimensional complex vector space. There are a group of representations (which are irreducible since is 1-dim) where for .
The character of such representation, that is the trace of for each and , is itself. In the case of representation of finite group, the character table of the group are rows of vectors such that each row is the character of one irreducible representation of , and these vectors form an orthonormal basis of the space of class functions that map from to by Schur's lemma. Now the group is no longer finite but still compact, and it preserves the orthonormality of character table. Each row of the table is the function of and the inner product between two class functions (all functions being class functions since is abelian) is defined as with the normalizing factor . The sequence is an orthonormal basis of the space of class functions .
For any representation of a finite group , can be expressed as the span ( are the irreps of ), such that . Similarly for and , . The Pontriagin dual is and for , is its Fourier transform for .
Gelfand transform
The Fourier transform is also a special case of Gelfand transform. In this particular context, it is closely related to the Pontryagin duality map defined above.
Given an abelian locally compact Hausdorff topological group , as before we consider space , defined using a Haar measure. With convolution as multiplication, is an abelian Banach algebra. It also has an involution * given by
Taking the completion with respect to the largest possibly -norm gives its enveloping -algebra, called the group -algebra of . (Any -norm on is bounded by the norm, therefore their supremum exists.)
Given any abelian -algebra , the Gelfand transform gives an isomorphism between and , where is the multiplicative linear functionals, i.e. one-dimensional representations, on with the weak-* topology. The map is simply given by
It turns out that the multiplicative linear functionals of , after suitable identification, are exactly the characters of , and the Gelfand transform, when restricted to the dense subset is the Fourier–Pontryagin transform.
Compact non-abelian groups
The Fourier transform can also be defined for functions on a non-abelian group, provided that the group is compact. Removing the assumption that the underlying group is abelian, irreducible unitary representations need not always be one-dimensional. This means the Fourier transform on a non-abelian group takes values as Hilbert space operators. The Fourier transform on compact groups is a major tool in representation theory and non-commutative harmonic analysis.
Let be a compact Hausdorff topological group. Let denote the collection of all isomorphism classes of finite-dimensional irreducible unitary representations, along with a definite choice of representation on the Hilbert space of finite dimension for each . If is a finite Borel measure on , then the Fourier–Stieltjes transform of is the operator on defined by
where is the complex-conjugate representation of acting on . If is absolutely continuous with respect to the left-invariant probability measure on , represented as
for some , one identifies the Fourier transform of with the Fourier–Stieltjes transform of .
The mapping
defines an isomorphism between the Banach space of finite Borel measures (see rca space) and a closed subspace of the Banach space consisting of all sequences indexed by of (bounded) linear operators for which the norm
is finite. The "convolution theorem" asserts that, furthermore, this isomorphism of Banach spaces is in fact an isometric isomorphism of C*-algebras into a subspace of . Multiplication on is given by convolution of measures and the involution * defined by
and has a natural -algebra structure as Hilbert space operators.
The Peter–Weyl theorem holds, and a version of the Fourier inversion formula (Plancherel's theorem) follows: if , then
where the summation is understood as convergent in the sense.
The generalization of the Fourier transform to the noncommutative situation has also in part contributed to the development of noncommutative geometry. In this context, a categorical generalization of the Fourier transform to noncommutative groups is Tannaka–Krein duality, which replaces the group of characters with the category of representations. However, this loses the connection with harmonic functions.
Alternatives
In signal processing terms, a function (of time) is a representation of a signal with perfect time resolution, but no frequency information, while the Fourier transform has perfect frequency resolution, but no time information: the magnitude of the Fourier transform at a point is how much frequency content there is, but location is only given by phase (argument of the Fourier transform at a point), and standing waves are not localized in time – a sine wave continues out to infinity, without decaying. This limits the usefulness of the Fourier transform for analyzing signals that are localized in time, notably transients, or any signal of finite extent.
As alternatives to the Fourier transform, in time–frequency analysis, one uses time–frequency transforms or time–frequency distributions to represent signals in a form that has some time information and some frequency information – by the uncertainty principle, there is a trade-off between these. These can be generalizations of the Fourier transform, such as the short-time Fourier transform, fractional Fourier transform, Synchrosqueezing Fourier transform, or other functions to represent signals, as in wavelet transforms and chirplet transforms, with the wavelet analog of the (continuous) Fourier transform being the continuous wavelet transform.
Example
The following figures provide a visual illustration of how the Fourier transform's integral measures whether a frequency is present in a particular function. The first image depicts the function which is a 3 Hz cosine wave (the first term) shaped by a Gaussian envelope function (the second term) that smoothly turns the wave on and off. The next 2 images show the product which must be integrated to calculate the Fourier transform at +3 Hz. The real part of the integrand has a non-negative average value, because the alternating signs of and oscillate at the same rate and in phase, whereas and oscillate at the same rate but with orthogonal phase. The absolute value of the Fourier transform at +3 Hz is 0.5, which is relatively large. When added to the Fourier transform at -3 Hz (which is identical because we started with a real signal), we find that the amplitude of the 3 Hz frequency component is 1.
However, when you try to measure a frequency that is not present, both the real and imaginary component of the integral vary rapidly between positive and negative values. For instance, the red curve is looking for 5 Hz. The absolute value of its integral is nearly zero, indicating that almost no 5 Hz component was in the signal. The general situation is usually more complicated than this, but heuristically this is how the Fourier transform measures how much of an individual frequency is present in a function
To re-enforce an earlier point, the reason for the response at Hz is because and are indistinguishable. The transform of would have just one response, whose amplitude is the integral of the smooth envelope: whereas is
Applications
Linear operations performed in one domain (time or frequency) have corresponding operations in the other domain, which are sometimes easier to perform. The operation of differentiation in the time domain corresponds to multiplication by the frequency, so some differential equations are easier to analyze in the frequency domain. Also, convolution in the time domain corresponds to ordinary multiplication in the frequency domain (see Convolution theorem). After performing the desired operations, transformation of the result can be made back to the time domain. Harmonic analysis is the systematic study of the relationship between the frequency and time domains, including the kinds of functions or operations that are "simpler" in one or the other, and has deep connections to many areas of modern mathematics.
Analysis of differential equations
Perhaps the most important use of the Fourier transformation is to solve partial differential equations.
Many of the equations of the mathematical physics of the nineteenth century can be treated this way. Fourier studied the heat equation, which in one dimension and in dimensionless units is
The example we will give, a slightly more difficult one, is the wave equation in one dimension,
As usual, the problem is not to find a solution: there are infinitely many. The problem is that of the so-called "boundary problem": find a solution which satisfies the "boundary conditions"
Here, and are given functions. For the heat equation, only one boundary condition can be required (usually the first one). But for the wave equation, there are still infinitely many solutions which satisfy the first boundary condition. But when one imposes both conditions, there is only one possible solution.
It is easier to find the Fourier transform of the solution than to find the solution directly. This is because the Fourier transformation takes differentiation into multiplication by the Fourier-dual variable, and so a partial differential equation applied to the original function is transformed into multiplication by polynomial functions of the dual variables applied to the transformed function. After is determined, we can apply the inverse Fourier transformation to find .
Fourier's method is as follows. First, note that any function of the forms
satisfies the wave equation. These are called the elementary solutions.
Second, note that therefore any integral
satisfies the wave equation for arbitrary . This integral may be interpreted as a continuous linear combination of solutions for the linear equation.
Now this resembles the formula for the Fourier synthesis of a function. In fact, this is the real inverse Fourier transform of and in the variable .
The third step is to examine how to find the specific unknown coefficient functions and that will lead to satisfying the boundary conditions. We are interested in the values of these solutions at . So we will set . Assuming that the conditions needed for Fourier inversion are satisfied, we can then find the Fourier sine and cosine transforms (in the variable ) of both sides and obtain
and
Similarly, taking the derivative of with respect to and then applying the Fourier sine and cosine transformations yields
and
These are four linear equations for the four unknowns and , in terms of the Fourier sine and cosine transforms of the boundary conditions, which are easily solved by elementary algebra, provided that these transforms can be found.
In summary, we chose a set of elementary solutions, parametrized by , of which the general solution would be a (continuous) linear combination in the form of an integral over the parameter . But this integral was in the form of a Fourier integral. The next step was to express the boundary conditions in terms of these integrals, and set them equal to the given functions and . But these expressions also took the form of a Fourier integral because of the properties of the Fourier transform of a derivative. The last step was to exploit Fourier inversion by applying the Fourier transformation to both sides, thus obtaining expressions for the coefficient functions and in terms of the given boundary conditions and .
From a higher point of view, Fourier's procedure can be reformulated more conceptually. Since there are two variables, we will use the Fourier transformation in both and rather than operate as Fourier did, who only transformed in the spatial variables. Note that must be considered in the sense of a distribution since is not going to be : as a wave, it will persist through time and thus is not a transient phenomenon. But it will be bounded and so its Fourier transform can be defined as a distribution. The operational properties of the Fourier transformation that are relevant to this equation are that it takes differentiation in to multiplication by and differentiation with respect to to multiplication by where is the frequency. Then the wave equation becomes an algebraic equation in :
This is equivalent to requiring unless . Right away, this explains why the choice of elementary solutions we made earlier worked so well: obviously will be solutions. Applying Fourier inversion to these delta functions, we obtain the elementary solutions we picked earlier. But from the higher point of view, one does not pick elementary solutions, but rather considers the space of all distributions which are supported on the (degenerate) conic .
We may as well consider the distributions supported on the conic that are given by distributions of one variable on the line plus distributions on the line as follows: if is any test function,
where , and , are distributions of one variable.
Then Fourier inversion gives, for the boundary conditions, something very similar to what we had more concretely above (put , which is clearly of polynomial growth):
and
Now, as before, applying the one-variable Fourier transformation in the variable to these functions of yields two equations in the two unknown distributions (which can be taken to be ordinary functions if the boundary conditions are or ).
From a calculational point of view, the drawback of course is that one must first calculate the Fourier transforms of the boundary conditions, then assemble the solution from these, and then calculate an inverse Fourier transform. Closed form formulas are rare, except when there is some geometric symmetry that can be exploited, and the numerical calculations are difficult because of the oscillatory nature of the integrals, which makes convergence slow and hard to estimate. For practical calculations, other methods are often used.
The twentieth century has seen the extension of these methods to all linear partial differential equations with polynomial coefficients, and by extending the notion of Fourier transformation to include Fourier integral operators, some non-linear equations as well.
Fourier-transform spectroscopy
The Fourier transform is also used in nuclear magnetic resonance (NMR) and in other kinds of spectroscopy, e.g. infrared (FTIR). In NMR an exponentially shaped free induction decay (FID) signal is acquired in the time domain and Fourier-transformed to a Lorentzian line-shape in the frequency domain. The Fourier transform is also used in magnetic resonance imaging (MRI) and mass spectrometry.
Quantum mechanics
The Fourier transform is useful in quantum mechanics in at least two different ways. To begin with, the basic conceptual structure of quantum mechanics postulates the existence of pairs of complementary variables, connected by the Heisenberg uncertainty principle. For example, in one dimension, the spatial variable of, say, a particle, can only be measured by the quantum mechanical "position operator" at the cost of losing information about the momentum of the particle. Therefore, the physical state of the particle can either be described by a function, called "the wave function", of or by a function of but not by a function of both variables. The variable is called the conjugate variable to . In classical mechanics, the physical state of a particle (existing in one dimension, for simplicity of exposition) would be given by assigning definite values to both and simultaneously. Thus, the set of all possible physical states is the two-dimensional real vector space with a -axis and a -axis called the phase space.
In contrast, quantum mechanics chooses a polarisation of this space in the sense that it picks a subspace of one-half the dimension, for example, the -axis alone, but instead of considering only points, takes the set of all complex-valued "wave functions" on this axis. Nevertheless, choosing the -axis is an equally valid polarisation, yielding a different representation of the set of possible physical states of the particle. Both representations of the wavefunction are related by a Fourier transform, such that
or, equivalently,
Physically realisable states are , and so by the Plancherel theorem, their Fourier transforms are also . (Note that since is in units of distance and is in units of momentum, the presence of the Planck constant in the exponent makes the exponent dimensionless, as it should be.)
Therefore, the Fourier transform can be used to pass from one way of representing the state of the particle, by a wave function of position, to another way of representing the state of the particle: by a wave function of momentum. Infinitely many different polarisations are possible, and all are equally valid. Being able to transform states from one representation to another by the Fourier transform is not only convenient but also the underlying reason of the Heisenberg uncertainty principle.
The other use of the Fourier transform in both quantum mechanics and quantum field theory is to solve the applicable wave equation. In non-relativistic quantum mechanics, Schrödinger's equation for a time-varying wave function in one-dimension, not subject to external forces, is
This is the same as the heat equation except for the presence of the imaginary unit . Fourier methods can be used to solve this equation.
In the presence of a potential, given by the potential energy function , the equation becomes
The "elementary solutions", as we referred to them above, are the so-called "stationary states" of the particle, and Fourier's algorithm, as described above, can still be used to solve the boundary value problem of the future evolution of given its values for . Neither of these approaches is of much practical use in quantum mechanics. Boundary value problems and the time-evolution of the wave function is not of much practical interest: it is the stationary states that are most important.
In relativistic quantum mechanics, Schrödinger's equation becomes a wave equation as was usual in classical physics, except that complex-valued waves are considered. A simple example, in the absence of interactions with other particles or fields, is the free one-dimensional Klein–Gordon–Schrödinger–Fock equation, this time in dimensionless units,
This is, from the mathematical point of view, the same as the wave equation of classical physics solved above (but with a complex-valued wave, which makes no difference in the methods). This is of great use in quantum field theory: each separate Fourier component of a wave can be treated as a separate harmonic oscillator and then quantized, a procedure known as "second quantization". Fourier methods have been adapted to also deal with non-trivial interactions.
Finally, the number operator of the quantum harmonic oscillator can be interpreted, for example via the Mehler kernel, as the generator of the Fourier transform .
Signal processing
The Fourier transform is used for the spectral analysis of time-series. The subject of statistical signal processing does not, however, usually apply the Fourier transformation to the signal itself. Even if a real signal is indeed transient, it has been found in practice advisable to model a signal by a function (or, alternatively, a stochastic process) which is stationary in the sense that its characteristic properties are constant over all time. The Fourier transform of such a function does not exist in the usual sense, and it has been found more useful for the analysis of signals to instead take the Fourier transform of its autocorrelation function.
The autocorrelation function of a function is defined by
This function is a function of the time-lag elapsing between the values of to be correlated.
For most functions that occur in practice, is a bounded even function of the time-lag and for typical noisy signals it turns out to be uniformly continuous with a maximum at .
The autocorrelation function, more properly called the autocovariance function unless it is normalized in some appropriate fashion, measures the strength of the correlation between the values of separated by a time lag. This is a way of searching for the correlation of with its own past. It is useful even for other statistical tasks besides the analysis of signals. For example, if represents the temperature at time , one expects a strong correlation with the temperature at a time lag of 24 hours.
It possesses a Fourier transform,
This Fourier transform is called the power spectral density function of . (Unless all periodic components are first filtered out from , this integral will diverge, but it is easy to filter out such periodicities.)
The power spectrum, as indicated by this density function , measures the amount of variance contributed to the data by the frequency . In electrical signals, the variance is proportional to the average power (energy per unit time), and so the power spectrum describes how much the different frequencies contribute to the average power of the signal. This process is called the spectral analysis of time-series and is analogous to the usual analysis of variance of data that is not a time-series (ANOVA).
Knowledge of which frequencies are "important" in this sense is crucial for the proper design of filters and for the proper evaluation of measuring apparatuses. It can also be useful for the scientific analysis of the phenomena responsible for producing the data.
The power spectrum of a signal can also be approximately measured directly by measuring the average power that remains in a signal after all the frequencies outside a narrow band have been filtered out.
Spectral analysis is carried out for visual signals as well. The power spectrum ignores all phase relations, which is good enough for many purposes, but for video signals other types of spectral analysis must also be employed, still using the Fourier transform as a tool.
Other notations
Other common notations for include:
In the sciences and engineering it is also common to make substitutions like these:
So the transform pair can become
A disadvantage of the capital letter notation is when expressing a transform such as or which become the more awkward and
In some contexts such as particle physics, the same symbol may be used for both for a function as well as it Fourier transform, with the two only distinguished by their argument I.e. would refer to the Fourier transform because of the momentum argument, while would refer to the original function because of the positional argument. Although tildes may be used as in to indicate Fourier transforms, tildes may also be used to indicate a modification of a quantity with a more Lorentz invariant form, such as , so care must be taken. Similarly, often denotes the Hilbert transform of .
The interpretation of the complex function may be aided by expressing it in polar coordinate form
in terms of the two real functions and where:
is the amplitude and
is the phase (see arg function).
Then the inverse transform can be written:
which is a recombination of all the frequency components of . Each component is a complex sinusoid of the form whose amplitude is and whose initial phase angle (at ) is .
The Fourier transform may be thought of as a mapping on function spaces. This mapping is here denoted and is used to denote the Fourier transform of the function . This mapping is linear, which means that can also be seen as a linear transformation on the function space and implies that the standard notation in linear algebra of applying a linear transformation to a vector (here the function ) can be used to write instead of . Since the result of applying the Fourier transform is again a function, we can be interested in the value of this function evaluated at the value for its variable, and this is denoted either as or as . Notice that in the former case, it is implicitly understood that is applied first to and then the resulting function is evaluated at , not the other way around.
In mathematics and various applied sciences, it is often necessary to distinguish between a function and the value of when its variable equals , denoted . This means that a notation like formally can be interpreted as the Fourier transform of the values of at . Despite this flaw, the previous notation appears frequently, often when a particular function or a function of a particular variable is to be transformed. For example,
is sometimes used to express that the Fourier transform of a rectangular function is a sinc function, or
is used to express the shift property of the Fourier transform.
Notice, that the last example is only correct under the assumption that the transformed function is a function of , not of .
As discussed above, the characteristic function of a random variable is the same as the Fourier–Stieltjes transform of its distribution measure, but in this context it is typical to take a different convention for the constants. Typically characteristic function is defined
As in the case of the "non-unitary angular frequency" convention above, the factor of 2 appears in neither the normalizing constant nor the exponent. Unlike any of the conventions appearing above, this convention takes the opposite sign in the exponent.
Computation methods
The appropriate computation method largely depends how the original mathematical function is represented and the desired form of the output function. In this section we consider both functions of a continuous variable, and functions of a discrete variable (i.e. ordered pairs of and values). For discrete-valued the transform integral becomes a summation of sinusoids, which is still a continuous function of frequency ( or ). When the sinusoids are harmonically related (i.e. when the -values are spaced at integer multiples of an interval), the transform is called discrete-time Fourier transform (DTFT).
Discrete Fourier transforms and fast Fourier transforms
Sampling the DTFT at equally-spaced values of frequency is the most common modern method of computation. Efficient procedures, depending on the frequency resolution needed, are described at . The discrete Fourier transform (DFT), used there, is usually computed by a fast Fourier transform (FFT) algorithm.
Analytic integration of closed-form functions
Tables of closed-form Fourier transforms, such as and , are created by mathematically evaluating the Fourier analysis integral (or summation) into another closed-form function of frequency ( or ). When mathematically possible, this provides a transform for a continuum of frequency values.
Many computer algebra systems such as Matlab and Mathematica that are capable of symbolic integration are capable of computing Fourier transforms analytically. For example, to compute the Fourier transform of one might enter the command into Wolfram Alpha.
Numerical integration of closed-form continuous functions
Discrete sampling of the Fourier transform can also be done by numerical integration of the definition at each value of frequency for which transform is desired. The numerical integration approach works on a much broader class of functions than the analytic approach.
Numerical integration of a series of ordered pairs
If the input function is a series of ordered pairs, numerical integration reduces to just a summation over the set of data pairs. The DTFT is a common subcase of this more general situation.
Tables of important Fourier transforms
The following tables record some closed-form Fourier transforms. For functions and denote their Fourier transforms by and . Only the three most common conventions are included. It may be useful to notice that entry 105 gives a relationship between the Fourier transform of a function and the original function, which can be seen as relating the Fourier transform and its inverse.
Functional relationships, one-dimensional
The Fourier transforms in this table may be found in or .
Square-integrable functions, one-dimensional
The Fourier transforms in this table may be found in , , or .
Distributions, one-dimensional
The Fourier transforms in this table may be found in or .
Two-dimensional functions
Formulas for general -dimensional functions
| Mathematics | Calculus and analysis | null |
52255 | https://en.wikipedia.org/wiki/Wild%20boar | Wild boar | The wild boar (Sus scrofa), also known as the wild swine, common wild pig, Eurasian wild pig, or simply wild pig, is a suid native to much of Eurasia and North Africa, and has been introduced to the Americas and Oceania. The species is now one of the widest-ranging mammals in the world, as well as the most widespread suiform. It has been assessed as least concern on the IUCN Red List due to its wide range, high numbers, and adaptability to a diversity of habitats. It has become an invasive species in part of its introduced range. Wild boars probably originated in Southeast Asia during the Early Pleistocene and outcompeted other suid species as they spread throughout the Old World.
, up to 16 subspecies are recognized, which are divided into four regional groupings based on skull height and lacrimal bone length. The species lives in matriarchal societies consisting of interrelated females and their young (both male and female). Fully grown males are usually solitary outside the breeding season. The wolf is the wild boar's main predator in most of its natural range except in the Far East and the Lesser Sunda Islands, where it is replaced by the tiger and Komodo dragon respectively. The wild boar has a long history of association with humans, having been the ancestor of most domestic pig breeds and a big-game animal for millennia. Boars have also re-hybridized in recent decades with feral pigs; these boar–pig hybrids have become a serious pest wild animal in the Americas and Australia.
Terminology
As true wild boars became extinct in Great Britain before the development of Modern English, the same terms are often used for both true wild boar and pigs, especially large or semi-wild ones. The English boar stems from the Old English , which is thought to be derived from the West Germanic , of unknown origin. Boar is sometimes used specifically to refer to males, and may also be used to refer to male domesticated pigs, especially breeding males that have not been castrated.
Sow, the traditional name for a female, again comes from Old English and Germanic; it stems from Proto-Indo-European, and is related to the Latin sus and Ancient Greek hus, and more closely to the New High German . The young may be called piglets or boarlets.
The animals' specific name scrofa is Latin for 'sow'.
Hunting
In hunting terminology, boars are given different designations according to their age:
Taxonomy and evolution
MtDNA studies indicate that the wild boar originated from islands in Southeast Asia such as Indonesia and the Philippines, and subsequently spread onto mainland Eurasia and North Africa. The earliest fossil finds of the species come from both Europe and Asia, and date back to the Early Pleistocene. By the late Villafranchian, S. scrofa largely displaced the related S. strozzii, a large, possibly swamp-adapted suid ancestral to the modern S. verrucosus throughout the Eurasian mainland, restricting it to insular Asia. Its closest wild relative is the bearded pig of Malacca and surrounding islands.
Subspecies
, 16 subspecies are recognised, which are divided into four regional groupings:
Western: Includes S. s. scrofa, S. s. meridionalis, S. s. algira, S. s. attila, S. s. lybicus, S. s. majori and S. s. nigripes. These subspecies are typically high-skulled (though lybicus and some scrofa are low-skulled), with thick underwool and (excepting scrofa and attila) poorly developed manes.
Indian: Includes S. s. davidi and S. s. cristatus. These subspecies have sparse or absent underwool, with long manes and prominent bands on the snout and mouth. While S. s. cristatus is high-skulled, S. s. davidi is low-skulled.
Eastern: Includes S. s. sibiricus, S. s. ussuricus, S. s. leucomystax, S. s. riukiuanus, S. s. taivanus and S. s. moupinensis. These subspecies are characterised by a whitish streak extending from the corners of the mouth to the lower jaw. With the exception of S. s. ussuricus, most are high-skulled. The underwool is thick, except in S. s. moupinensis, and the mane is largely absent.
Indonesian: Represented solely by S. s. vittatus, it is characterised by its sparse body hair, lack of underwool, fairly long mane, a broad reddish band extending from the muzzle to the sides of the neck. It is the most basal of the four groups, having the smallest relative brain size, more primitive dentition and unspecialised cranial structure.
Domestication
With the exception of domestic pigs in Timor and Papua New Guinea (which appear to be of Sulawesi warty pig stock), the wild boar is the ancestor of most pig breeds. Archaeological evidence suggests that pigs were domesticated from wild boar as early as 13,000–12,700 BCE in the Near East in the Tigris Basin, being managed in the wild in a way similar to the way they are managed by some modern New Guineans. Remains of pigs have been dated to earlier than 11,400 BCE in Cyprus. Those animals must have been introduced from the mainland, which suggests domestication in the adjacent mainland by then. There was also a separate domestication in China, which took place about 8,000 years ago.
DNA evidence from sub-fossil remains of teeth and jawbones of Neolithic pigs shows that the first domestic pigs in Europe had been brought from the Near East. This stimulated the domestication of local European wild boars, resulting in a third domestication event with the Near Eastern genes dying out in European pig stock. Modern domesticated pigs have involved complex exchanges, with European domesticated lines being exported in turn to the ancient Near East. Historical records indicate that Asian pigs were introduced into Europe during the 18th and early 19th centuries. Domestic pigs tend to have much more developed hindquarters than their wild boar ancestors, to the point where 70% of their body weight is concentrated in the posterior, which is the opposite of wild boar, where most of the muscles are concentrated on the head and shoulders.
Synonymous species
The Heude's pig (Sus bucculentus), also known as the Indochinese warty pig or Vietnam warty pig, was an alleged pig species found in Laos and Vietnam. It was virtually unknown and was feared extinct, until the discovery of a skull from a recently killed individual in the Annamite Range, Laos, in 1995. Subsequent studies indicated that Sus bucculentus was not a valid taxon. As of 2022 the Mammal Diversity Database included it in Sus scrofa.
Description
The wild boar is a bulky, massively built suid with short and relatively thin legs. The trunk is short and robust, while the hindquarters are comparatively underdeveloped. The region behind the shoulder blades rises into a hump and the neck is short and thick to the point of being nearly immobile. The animal's head is very large, taking up to one-third of the body's entire length. The structure of the head is well suited for digging. The head acts as a plough, while the powerful neck muscles allow the animal to upturn considerable amounts of soil: it is capable of digging into frozen ground and can upturn rocks weighing . The eyes are small and deep-set and the ears long and broad. The species has well developed canine teeth, which protrude from the mouths of adult males. The medial hooves are larger and more elongated than the lateral ones and are capable of quick movements. The animal can run at a maximum speed of 40 km/h (25 mph) and jump at a height of .
Sexual dimorphism is very pronounced in the species, with males being typically 5–10% larger and 20–30% heavier than females. Males also sport a mane running down the back, which is particularly apparent during autumn and winter. The canine teeth are also much more prominent in males and grow throughout life. The upper canines are relatively short and grow sideways early in life, though they gradually curve upwards. The lower canines are much sharper and longer, with the exposed parts measuring in length. In the breeding period, males develop a coating of subcutaneous tissue, which may be thick, extending from the shoulder blades to the rump, thus protecting vital organs during fights. Males sport a roughly chicken egg–sized sac of unclear function near the opening of the penis that collects urine and emits a sharp odour.
Adult size and weight is largely determined by environmental factors; boars living in arid areas with little productivity tend to attain smaller sizes than their counterparts inhabiting areas with abundant food and water. In most of Europe, males average in weight, in shoulder height and in body length, whereas females average in weight, in shoulder height and in body length. In Europe's Mediterranean regions, males may reach average weights as low as and females , with shoulder heights of . In the more productive areas of Eastern Europe, males average in weight, in shoulder height and in body length, while females weigh , reach in shoulder height, and reach in body length. In Western and Central Europe, the largest males weigh and females . In Northeastern Asia, large males can reach brown bear-like sizes, weighing and measuring in shoulder height. Some adult males in Ussuriland and Manchuria have been recorded to weigh and measure in shoulder height. Adults of this size are generally immune from wolf predation. Such giants are rare in modern times, as past overhunting has prevented animals from attaining their full growth.
The winter coat consists of long, coarse bristles underlaid with short brown downy fur. The length of these bristles varies along the body, with the shortest being around the face and limbs and the longest running along the back. These back bristles form the aforementioned mane prominent in males and stand erect when the animal is agitated. Colour is highly variable; specimens around Lake Balkhash are very lightly coloured, and can even be white, while some boars from Belarus and Ussuriland can be black. Some subspecies sport a light-coloured patch running backward from the corners of the mouth. Coat colour also varies with age, with piglets having light brown or rusty-brown fur with pale bands extending from the flanks and back.
The wild boar produces a number of different sounds which are divided into three categories:
Contact calls: Grunting noises which differ in intensity according to the situation. Adult males are usually silent, while females frequently grunt and piglets whine. When feeding, boars express their contentment through purring. Studies have shown that piglets imitate the sounds of their mother, thus different litters may have unique vocalisations.
Alarm calls: Warning cries emitted in response to threats. When frightened, boars make loud huffing sounds or emit screeches transcribed as .
Combat calls: High-pitched, piercing cries.
Its sense of smell is very well developed to the point that the animal is used for drug detection in Germany. Its hearing is also acute, though its eyesight is comparatively weak, lacking color vision and being unable to recognise a standing human away.
Pigs are one of four known mammalian taxa which possess mutations in the nicotinic acetylcholine receptor that protect against snake venom. Mongooses, honey badgers, hedgehogs, and pigs all have modifications to the receptor pocket which prevents the snake venom α-neurotoxin from binding. These represent four separate, independent mutations.
Social behaviour and life cycle
Boars are typically social animals, living in female-dominated sounders consisting of barren sows and mothers with young led by an old matriarch. Male boars leave their sounder at the age of 8–15 months, while females either remain with their mothers or establish new territories nearby. Subadult males may live in loosely knit groups, while adult and elderly males tend to be solitary outside the breeding season.
The period in most areas lasts from November to January, though most mating only lasts a month and a half. Prior to mating, the males develop their subcutaneous armour in preparation for confronting rivals. The testicles double in size and the glands secrete a foamy yellowish liquid. Once ready to reproduce, males travel long distances in search of a sounder of sows, eating little on the way. Once a sounder has been located, the male drives off all young animals and persistently chases the sows. At this point, the male fiercely fights potential rivals. A single male can mate with 5–10 sows. By the end of the rut, males are often badly mauled and have lost 20% of their body weight, with bite-induced injuries to the penis being common. The gestation period varies according to the age of the expecting mother. For first-time breeders, it lasts 114–130 days, while it lasts 133–140 days in older sows. Farrowing occurs between March and May, with litter sizes depending on the age and nutrition of the mother. The average litter consists of 4–6 piglets, with the maximum being 10–12. The piglets are whelped in a nest constructed from twigs, grasses and leaves. Should the mother die prematurely, the piglets are adopted by the other sows in the sounder.
Newborn piglets weigh around 600–1,000 grams, lacking underfur and bearing a single milk incisor and canine on each half of the jaw. There is intense competition between the piglets over the most milk-rich nipples, as the best-fed young grow faster and have stronger constitutions. The piglets do not leave the lair for their first week of life. Should the mother be absent, the piglets lie closely pressed to each other. By two weeks of age, the piglets begin accompanying their mother on her journeys. Should danger be detected, the piglets take cover or stand immobile, relying on their camouflage to keep them hidden. The neonatal coat fades after three months, with adult colouration being attained at eight months. Although the lactation period lasts 2.5–3.5 months, the piglets begin displaying adult feeding behaviours at the age of 2–3 weeks. The permanent dentition is fully formed by 1–2 years. With the exception of the canines in males, the teeth stop growing during the middle of the fourth year. The canines in old males continue to grow throughout their lives, curving strongly as they age. Sows attain sexual maturity at the age of one year, with males attaining it a year later. However, estrus usually first occurs after two years in sows, while males begin participating in the rut after 4–5 years, as they are not permitted to mate by the older males. The maximum lifespan in the wild is 10–14 years, though few specimens survive past 4–5 years. Boars in captivity have lived for 20 years.
Behaviour and ecology
Habitat and sheltering
The wild boar inhabits a diverse array of habitats from boreal taigas to deserts. In mountainous regions, it can even occupy alpine zones, occurring up to in the Carpathians, in the Caucasus and up to in the mountains in Central Asia and Kazakhstan. In order to survive in a given area, wild boars require a habitat fulfilling three conditions: heavily brushed areas providing shelter from predators, water for drinking and bathing purposes and an absence of regular snowfall.
The main habitats favored by boars in Europe are deciduous and mixed forests, with the most favorable areas consisting of forest composed of oak and beech enclosing marshes and meadows. In the Białowieża Forest, the animal's primary habitat consists of well-developed broad-leaved and mixed forests, along with marshy mixed forests, with coniferous forests and undergrowths being of secondary importance. Forests made up entirely of oak groves and beeches are used only during the fruit-bearing season. This is in contrast to the Caucasian and Transcaucasian mountain areas, where boars will occupy such fruit-bearing forests year-round. In the mountainous areas of the Russian Far East, the species inhabits nutpine groves, hilly mixed forests where Mongolian oak and Korean pine are present, swampy mixed taiga and coastal oak forests. In Transbaikalia, boars are restricted to river valleys with nut pine and shrubs. Boars are regularly encountered in pistachio groves in winter in some areas of Tajikistan and Turkmenistan, while in spring they migrate to open deserts; boar have also colonized deserts in several areas they have been introduced to.
On the islands of Komodo and Rinca, the boar mostly inhabits savanna or open monsoon forests, avoiding heavily forested areas unless pursued by humans. Wild boar are known to be competent swimmers, capable of covering long distances. In 2013, one boar was reported to have completed the swim from France to Alderney in the Channel Islands. Due to concerns about disease, it was shot and incinerated.
Wild boar rest in shelters, which contain insulating material like spruce branches and dry hay. These resting places are occupied by whole families (though males lie separately) and are often located in the vicinity of streams, in swamp forests and in tall grass or shrub thickets. Boars never defecate in their shelters and will cover themselves with soil and pine needles when irritated by insects.
Diet
The wild boar is a highly versatile omnivore, whose diversity in choice of food is comparable to that of humans. Their foods can be divided into four categories:
Rhizomes, roots, tubers and bulbs, all of which are dug up throughout the year in the animal's whole range.
Nuts, berries and seeds, which are consumed when ripened and are dug up from the snow when necessary.
Leaves, bark, twigs and shoots, along with garbage.
Earthworms, insects, mollusks, fish, rodents, insectivores, bird eggs, lizards, snakes, frogs and carrion. Most of these prey items are taken in warm periods.
A boar needs around 4,000–4,500 calories of food per day, though this required amount increases during winter and pregnancy, with the majority of its diet consisting of food items dug from the ground, like underground plant material and burrowing animals. Acorns and beechnuts are invariably its most important food items in temperate zones, as they are rich in the carbohydrates necessary for the buildup of fat reserves needed to survive lean periods. In Western Europe, underground plant material favoured by boars includes bracken, willow herb, bulbs, meadow herb roots and bulbs and the bulbs of cultivated crops. Such food is favoured in early spring and summer, but may also be eaten in autumn and winter during beechnut and acorn crop failures. Should regular wild foods become scarce, boars will eat tree bark and fungi, as well as visit cultivated potato and artichoke fields. Boar soil disturbance and foraging have been shown to facilitate invasive plants. Boars of the vittatus subspecies in Ujung Kulon National Park in Java differ from most other populations by their primarily frugivorous diet, which consists of 50 different fruit species, especially figs, thus making them important seed dispersers. The wild boar can consume numerous genera of poisonous plants without ill effect, including Aconitum, Anemone, Calla, Caltha, Ferula and Pteridium.
Boars may occasionally prey on small vertebrates like newborn deer fawns, leporids and galliform chicks. Boars inhabiting the Volga Delta and near some lakes and rivers of Kazakhstan have been recorded to feed extensively on fish like carp and Caspian roach. Boars in the former area will also feed on cormorant and heron chicks, bivalved molluscs, trapped muskrats and mice. There is at least one record of a boar killing and eating a bonnet macaque in southern India's Bandipur National Park, though this may have been a case of intraguild predation, brought on by interspecific competition for human handouts. There is also at least one recorded case of a group of wild boar attacking, killing, and eating an adult, healthy female axis deer (Axis axis) as a pack.
Stable isotope analysis of fossil wild boar tooth enamel from the late Middle Pleistocene found in Thailand indicate that it fed on a versatile mixed vegetation.
Predators
Piglets are vulnerable to attack from medium-sized felids like Eurasian lynx (Lynx lynx), jungle cats (Felis chaus), and snow leopards (Panthera uncia), as well as other carnivorans like brown bears (Ursus arctos) and yellow-throated martens (Martes flavigula).
The wolf (Canis lupus) is the main predator of wild boar throughout most of its range. A single wolf can kill around 50 to 80 boars of differing ages in one year. In Italy and Belarus' Belovezhskaya Pushcha National Park, boars are the wolf's primary prey, despite an abundance of alternative, less powerful ungulates. Wolves are particularly threatening during the winter, when deep snow impedes the boars' movements. In the Baltic regions, heavy snowfall can allow wolves to eliminate boars from an area almost completely. Wolves primarily target piglets and subadults and only rarely attack adult sows. Adult males are usually avoided entirely. Dholes (Cuon alpinus) may also prey on boars, to the point of keeping their numbers down in northwestern Bhutan, despite there being many more cattle in the area.
Leopards (Panthera pardus) are predators of wild boar in the Caucasus (particularly Transcaucasia), the Russian Far East, India, China and Iran. In most areas, boars constitute only a small part of the leopard's diet. However, in Iran's Sarigol National Park, boars are the second most frequently targeted prey species after mouflon (Ovis gmelini), though adult individuals are generally avoided, as they are above the leopard's preferred weight range of . This dependence on wild boar is largely due in part to the local leopard subspecies' large size.
Boars of all ages were once the primary prey of the tiger (Panthera tigris) in Transcaucasia, Kazakhstan, Middle Asia and the Far East up until the late 19th century. In modern times, tiger numbers are too low to have a limiting effect on boar populations. A single tiger can systematically destroy an entire sounder by preying on its members one by one, before moving on to another sounder. Tigers have been noted to chase boars for longer distances than with other prey. In two rare cases, boars were reported to gore a small tiger and a tigress to death in self-defense. A "large male tiger" died of wounds inflicted by an old wild boar it had killed in "a battle royal" between the two animals.
In the Amur region, wild boars are one of the two most important prey species for Siberian tigers, alongside the Manchurian wapiti (Cervus canadensis xanthopygus), with the two species collectively comprising roughly 80% of the felid's prey. In Sikhote Alin, a tiger can kill 30–34 boars a year. Studies of tigers in India indicate that boars are usually secondary in preference to various cervids and bovids, though when boars are targeted, healthy adults are caught more frequently than young and sick specimens.
On the islands of Komodo, Rinca and Flores, the boar's main predator is the Komodo dragon (Varanus komodoensis).
Distribution and habitat
Reconstructed range
The species originally occurred in North Africa and much of Eurasia; from the British Isles to Korea and the Sunda Islands. The northern limit of its range extended from southern Scandinavia to southern Siberia and Japan. Within this range, it was only absent in extremely dry deserts and alpine zones. It was once found in North Africa along the Nile valley up to Khartoum and north of the Sahara. The species occurs on a few Ionian and Aegean Islands, sometimes swimming between islands. The reconstructed northern boundary of the animal's Asian range ran from Lake Ladoga (at 60°N) through the area of Novgorod and Moscow into the southern Urals, where it reached 52°N. From there, the boundary passed Ishim and farther east the Irtysh at 56°N. In the eastern Baraba steppe (near Novosibirsk) the boundary turned steep south, encircled the Altai Mountains and went again eastward including the Tannu-Ola Mountains and Lake Baikal. From here, the boundary went slightly north of the Amur River eastward to its lower reaches at the Sea of Okhotsk. On Sakhalin, there are only fossil reports of wild boar. The southern boundaries in Europe and Asia were almost invariably identical to the seashores of these continents. It is absent in the dry regions of Mongolia from 44 to 46°N southward, in China westward of Sichuan and in India north of the Himalayas. It is absent in the higher elevations of the Pamir and the Tian Shan, though they do occur in the Tarim Basin and on the lower slopes of the Tien Shan.
Present range
In recent centuries, the range of wild boar has changed dramatically, largely due to hunting by humans and more recently because of captive wild boar escaping into the wild. Prior to the 20th century, boar populations had declined in numerous areas, with British populations probably becoming extinct during the 13th century. In the warm period after the ice age, wild boar lived in the southern parts of Sweden and Norway and north of Lake Ladoga in Karelia. It was previously thought that the species did not live in Finland during prehistory because no prehistoric wild boar bones had been found within the borders of the country. It was not until 2013, when a wild boar bone was found in Askola, that the species was found to have lived in Finland more than 8,000 years ago. It is believed, however, that man prevented its establishment by hunting. In Denmark, the last boar was shot at the beginning of the 19th century, and by 1900 they were absent in Tunisia and Sudan and large areas of Germany, Austria and Italy. In Russia, they were extirpated in wide areas by the 1930s. The last boar in Egypt reportedly died on 20 December 1912 in the Giza Zoo, with wild populations having disappeared by 1894–1902. Prince Kamal el Dine Hussein attempted to repopulate Wadi El Natrun with boars of Hungarian stock, but they were quickly exterminated by poachers.
A revival of boar populations began in the middle of the 20th century. By 1950, wild boar had once again reached their original northern boundary in many parts of their Asiatic range. By 1960, they reached Leningrad and Moscow and by 1975, they were to be found in Archangelsk and Astrakhan. In the 1970s they again occurred in Denmark and Sweden, where captive animals escaped and now survive in the wild. In England, wild boar populations re-established themselves in the 1990s, after escaping from specialist farms that had imported European stock.
Status in Great Britain
By the 11th century, wild boars were apparently already becoming rare in Britain. A 1087 forestry law enacted by William the Conqueror punished through blinding the unlawful killing of a boar. Charles I attempted to reintroduce the species into the New Forest, but this population was exterminated in the 17th century during the English Civil War. Between their medieval extinction and the 1980s, when wild boar farming began, only a handful of captive wild boar, imported from the continent, were present in Britain. Occasional escapes of wild boar from wildlife parks have occurred as early as the 1970s, but since the early 1990s significant populations have re-established themselves after escapes from farms, the number of which has increased as the demand for meat from the species has grown. A 1998 MAFF (now DEFRA) study on wild boar living wild in Britain confirmed the presence of two populations of wild boar living in Britain; one in Kent/East Sussex and another in Dorset.
Another DEFRA report, in February 2008, confirmed the existence of these two sites as 'established breeding areas' and identified a third in Gloucestershire/Herefordshire; in the Forest of Dean/Ross on Wye area. A 'new breeding population' was also identified in Devon. There is another significant population in Dumfries and Galloway. Populations estimates were as follows:
The largest population, in Kent/East Sussex, was then estimated at 200 animals in the core distribution area.
The smallest, in west Dorset, was estimated to be fewer than 50 animals.
Since winter 2005–2006 significant escapes/releases have also resulted in animals colonizing areas around the fringes of Dartmoor, in Devon. These are considered as an additional single 'new breeding population' and currently estimated to be up to 100 animals.
Population estimates for the Forest of Dean are disputed as, at the time that the DEFRA population estimate was 100, a photo of a boar sounder in the forest near Staunton with over 33 animals visible was published and at about the same time over 30 boar were seen in a field near the original escape location of Weston under Penyard many kilometres or miles away. In early 2010 the Forestry Commission embarked on a cull, with the aim of reducing the boar population from an estimated 150 animals to 100. By August it was stated that efforts were being made to reduce the population from 200 to 90, but that only 25 had been killed. The failure to meet cull targets was confirmed in February 2011.
Wild boars have crossed the River Wye into Monmouthshire, Wales. Iolo Williams, the BBC Wales wildlife expert, attempted to film Welsh boar in late 2012. Many other sightings, across the UK, have also been reported. The effects of wild boar on the U.K.'s woodlands were discussed with Ralph Harmer of the Forestry Commission on the 's Farming Today radio programme in 2011. The programme prompted activist writer George Monbiot to propose a thorough population study, followed by the introduction of permit-controlled culling.
In Scotland, wild boar are professionally referred to as 'feral pigs' as the genetics of the established feral populations may come from a mix of both wild boar and domestic pigs. They are now known to be present in Dumfries and Galloway and a number of sites in the Highlands, mainly centred around the Loch Ness area. They can be killed there legally without a license and are culled by land managers as wild populations appear occasionally. , an agency that advises the Scottish government estimates that Scotland is home to a few thousand wild boars. Locals around Loch Ness that were interviewed by The New York Times believed that the boars were becoming an increasingly worse problem, and farmers noted that they had killed and eaten several lambs.
Introduction to North America
Wild boars are an invasive species in the Americas, having been introduced by European explorers and settlers in the 16th century to serve as a source of food. Wild boars now cause problems including out-competing native species for food, destroying the nests of ground-nesting species, killing fawns and young domestic livestock, destroying agricultural crops, eating tree seeds and seedlings, destroying native vegetation and wetlands through wallowing, damaging water quality, coming into violent conflict with humans and pets and carrying pig and human diseases including brucellosis, trichinosis and pseudorabies. In some jurisdictions, it is illegal to import, breed, release, possess, sell, distribute, trade, transport, hunt, or trap Eurasian boars. Hunting and trapping is done systematically, to increase the chance of eradication and to remove the incentive to illegally release boars, which have mostly been spread deliberately by sport hunters.
History
While domestic pigs, both captive and feral (popularly termed "razorbacks"), have been in North America since the earliest days of European colonization, pure wild boars were not introduced into the New World until the 19th century. The suids were released into the wild by wealthy landowners as big game animals. The initial introductions took place in fenced enclosures, though several escapes occurred, with the escapees sometimes intermixing with already established feral pig populations.
The first of these introductions occurred in New Hampshire in 1890. Thirteen wild boars from Germany were purchased by Austin Corbin from Carl Hagenbeck and released into a game preserve in Sullivan County. Several of these boars escaped, though they were quickly hunted down by locals. Two further introductions were made from the original stocking, with several escapes taking place due to breaches in the game preserve's fencing. These escapees have ranged widely, with some specimens having been observed crossing into Vermont.
In 1902, 15–20 wild boar from Germany were released into a estate in Hamilton County, New York. Several specimens escaped six years later, dispersing into the William C. Whitney Wilderness Area, with their descendants surviving for at least 20 years.
The most extensive boar introduction in the US took place in western North Carolina in 1912, when 13 boars of undetermined European origin were released into two fenced enclosures in a game preserve in Hooper Bald, Graham County. Most of the specimens remained in the preserve for the next decade, until a large-scale hunt caused the remaining animals to break through their confines and escape. Some of the boars migrated to Tennessee, where they intermixed with both free-ranging and feral pigs in the area. In 1924, a dozen Hooper Bald wild pigs were shipped to California and released in a property between Carmel Valley and the Los Padres National Forest. These hybrid boar were later used as breeding stock on various private and public lands throughout the state, as well as in other states like Florida, Georgia, South Carolina, West Virginia and Mississippi.
Several wild boars from Leon Springs and the San Antonio, Saint Louis and San Diego Zoos were released in the Powder Horn Ranch in Calhoun County, Texas, in 1939. These specimens escaped and established themselves in surrounding ranchlands and coastal areas, with some crossing the Espiritu Santo Bay and colonizing Matagorda Island. Descendants of the Powder Horn Ranch boars were later released onto San José Island and the coast of Chalmette, Louisiana.
Wild boar of unknown origin were stocked in a ranch in the Edwards Plateau in the 1940s, only to escape during a storm and hybridize with local feral pig populations, later spreading into neighboring counties.
Starting in the mid-1980s, several boars purchased from the San Diego Zoo and Tierpark Berlin were released into the United States. A decade later, more specimens from farms in Canada and Białowieża Forest were let loose. In recent years, wild pig populations have been reported in 44 states within the US, most of which are likely wild boar–feral hog hybrids. Pure wild boar populations may still be present, but are extremely localized.
Introduction and lack of control in South America
In South America, the European boar is believed to have been introduced for the first time in Argentina and Uruguay around the 20th century for breeding purposes. In Brazil, the creation of wild boar and hybrids started on a large scale in the mid-1990s. With the invasion of wild boar that crossed the border and entered Rio Grande do Sul around 1989, and the escape and intentional release by several Brazilian breeders in the late 1990s – in response to a IBAMA decision against the import and breeding of wild boar in 1998 – numerous feral species formed a growing population, which progressively advances in Brazilian territory.
Pest control in Brazil
As a form of control for the wild boar population, hunting and killing are allowed for Collectors, Shooters and Hunters (CACs) duly registered by the environmental control agency, IBAMA, which, on the other hand, seeks to encourage the preservation of similar species of native peccaries, such as the "queixada" and the "caititu".
Effect on other habitats
Wild boars negatively impact other habitats through the destruction of the environment, or homes of wildlife. When wild boars invade new areas, they adapt to the new area by trampling and rooting, as well as displacing many saplings/nutrients. This causes a decrease in growing of many plants and trees. Water is also affected negatively by wild boars. When wild boars are active in streams, or small pools of water, it causes increased turbidity (excessive silt and particle suspension). In some cases, the fecal coliform concentration increases to dangerous levels because of wild boars. Aquatic wildlife is affected, more prominently fish, and amphibians. Wild boars have caused a great decrease in over 300 animal or plant species, 250 being endangered or threatened.
The boars cause many habitats to become less diverse because of their feeding behaviors and predation. Wild boars will dig up eggs of species and eat them, as well as killing other wildlife for food. When these boars compete with other species for resources, they usually come out successful. A study published in the Journal of Experimental Marine Biology and Ecology was conducted on the results of Feral Swine control. Only two years after the control started, the amount of turtle nests jumped from 57 to 143, and the turtle nest predation percent dropped from 74 to 15. They kill and eat deers, lizards, birds, snakes, and more. These boars are called "opportunist omnivores", which means they eat almost anything. This means they can survive almost anywhere. A big surplus of food and the ability to adapt to any new place causes lots of breeding. All of these factors make it difficult to get rid of wild boars.
Diseases and parasites
Wild boars are known to host at least 20 different parasitic worm species, with maximum infections occurring in summer. Young animals are vulnerable to helminths like Metastrongylus, which are consumed by boars through earthworms and cause death by parasitising the lungs. Wild boar also carry parasites known to infect humans, including Gastrodiscoides, Trichinella spiralis, Taenia solium, Balantidium coli and Toxoplasma gondii. Wild boar in southern regions are frequently infested with ticks (Dermacentor, Rhipicephalus, and Hyalomma) and hog lice. The species also suffers from blood-sucking flies, which it escapes by bathing frequently or hiding in dense shrubs.
Swine plague spreads very quickly in wild boar, with epizootics being recorded in Germany, Poland, Hungary, Belarus, the Caucasus, the Far East, Kazakhstan and other regions. Foot-and-mouth disease can also take on epidemic proportions in boar populations. The species occasionally, but rarely contracts Pasteurellosis, hemorrhagic sepsis, tularemia, and anthrax. Wild boar may on occasion contract swine erysipelas through rodents or hog lice and ticks.
Relationships with humans
In culture
The wild boar features prominently in the cultures of Indo-European people, many of which saw the animal as embodying warrior virtues. Cultures throughout Europe and Asia Minor saw the killing of a boar as proof of one's valor and strength. Neolithic hunter gatherers depicted reliefs of ferocious wild boars on their temple pillars at Göbekli Tepe some 11,600 years ago. Virtually all heroes in Greek mythology fight or kill a boar at one point. The demigod Herakles' third labour involves the capture of the Erymanthian Boar, Theseus slays the wild sow Phaea, and a disguised Odysseus is recognised by his handmaiden Eurycleia by the scars inflicted on him by a boar during a hunt in his youth. To the mythical Hyperboreans, the boar represented spiritual authority. Several Greek myths use the boar as a symbol of darkness, death and winter. One example is the story of the youthful Adonis, who is killed by a boar and is permitted by Zeus to depart from Hades only during the spring and summer period. This theme also occurs in Irish and Egyptian mythology, where the animal is explicitly linked to the month of October, therefore autumn. This association likely arose from aspects of the boar's actual nature. Its dark colour was linked to the night, while its solitary habits, proclivity to consume crops and nocturnal nature were associated with evil. The foundation myth of Ephesus has the city being built over the site where Prince Androklos of Athens killed a boar. Boars were frequently depicted on Greek funerary monuments alongside lions, representing gallant losers who have finally met their match, as opposed to victorious hunters as lions are. The theme of the doomed, yet valorous boar warrior also occurred in Hittite culture, where it was traditional to sacrifice a boar alongside a dog and a prisoner of war after a military defeat.
The boar as a warrior also appears in Germanic cultures, with its image having been frequently engraved on shields and swords. They also feature on Germanic boar helmets, such as the Benty Grange helmet, where it was believed to offer protection to the wearer and has been theorised to have been used in spiritual transformations into swine, similar to berserkers. The boar features heavily in religious practice in Germanic paganism where it is closely associated with Freyr and has also been suggested to have been a totemic animal to the Swedes, especially to the Yngling royal dynasty who claimed descent from the god.
According to Tacitus, the Baltic Aesti featured boars on their helmets and may have also worn boar masks. The boar and pig were held in particularly high esteem by the Celts, who considered them to be their most important sacred animal. Some Celtic deities linked to boars include Moccus and Veteris. It has been suggested that some early myths surrounding the Welsh hero Culhwch involved the character being the son of a boar god. Nevertheless, the importance of the boar as a culinary item among Celtic tribes may have been exaggerated in popular culture by the Asterix series, as wild boar bones are rare among Celtic archaeological sites and the few that do occur show no signs of butchery, having probably been used in sacrificial rituals.
The boar also appears in Vedic mythology and Hindu mythology. A story present in the Brahmanas has the god Indra slaying an avaricious boar, who has stolen the treasure of the asuras, then giving its carcass to the god Vishnu, who offered it as a sacrifice to the gods. In the story's retelling in the Charaka Samhita, the boar is described as a form of Prajapati and is credited with having raised the Earth from the primeval waters. In the Ramayana and the Puranas, the same boar is portrayed as Varaha, an avatar of Vishnu.
In Japanese culture, the boar is widely seen as a fearsome and reckless animal, to the point that several words and expressions in Japanese referring to recklessness include references to boars. The boar is the last animal of the Oriental zodiac, with people born during the year of the Pig being said to embody the boar-like traits of determination and impetuosity. Among Japanese hunters, the boar's courage and defiance is a source of admiration and it is not uncommon for hunters and mountain people to name their sons after the animal inoshishi (猪). Boars are also seen as symbols of fertility and prosperity; in some regions, it is thought that boars are drawn to fields owned by families including pregnant women, and hunters with pregnant wives are thought to have greater chances of success when boar hunting. The animal's link to prosperity was illustrated by its inclusion on the ¥10 note during the Meiji period and it was once believed that a man could become wealthy by keeping a clump of boar hair in his wallet.
In the folklore of the Mongol Altai Uriankhai tribe, the wild boar was associated with the watery underworld, as it was thought that the spirits of the dead entered the animal's head, to be ultimately transported to the water. Prior to the conversion to Islam, the Kyrgyz people believed that they were descended from boars and thus did not eat pork. In Buryat mythology, the forefathers of the Buryats descended from heaven and were nourished by a boar. In China, the boar is the emblem of the Miao people.
The boar (sanglier) is frequently displayed in English, Scottish and Welsh heraldry. As with the lion, the boar is often shown as armed and langued. As with the bear, Scottish and Welsh heraldry displays the boar's head with the neck cropped, unlike the English version, which retains the neck. The white boar served as the badge of King Richard III of England, who distributed it among his northern retainers during his tenure as Duke of Gloucester.
As a game animal and food source
Humans have been hunting boar for millennia, the earliest artistic depictions of such activities dating back to the Upper Paleolithic. The animal was seen as a source of food among the Ancient Greeks, as well as a sporting challenge and source of epic narratives. The Romans inherited this tradition, with one of its first practitioners being Scipio Aemilianus. Boar hunting became particularly popular among the young nobility during the 3rd century BC as preparation for manhood and battle. A typical Roman boar hunting tactic involved surrounding a given area with large nets, then flushing the boar with dogs and immobilizing it with smaller nets. The animal would then be dispatched with a venabulum, a short spear with a crossguard at the base of the blade. More than their Greek predecessors, the Romans extensively took inspiration from boar hunting in their art and sculpture. With the ascension of Constantine the Great, boar hunting took on Christian allegorical themes, with the animal being portrayed as a "black beast" analogous to the dragon of Saint George.
Boar hunting continued after the fall of the Western Roman Empire, though the Germanic tribes considered the red deer to be a more noble and worthy quarry. The post-Roman nobility hunted boar as their predecessors did, but primarily as training for battle rather than sport. It was not uncommon for medieval hunters to deliberately hunt boars during the breeding season when the animals were more aggressive. During the Renaissance, when deforestation and the introduction of firearms reduced boar numbers, boar hunting became the sole prerogative of the nobility, one of many charges brought up against the rich during the German Peasants' War and the French Revolution.
During the mid-20th century, 7,000–8,000 boars were caught in the Caucasus, 6,000–7,000 in Kazakhstan and about 5,000 in Central Asia during the Soviet period, primarily through the use of dogs and beats. In Nepal, farmers and poachers eliminate boars by baiting balls of wheat flour containing explosives with kerosene oil, with the animals' chewing motions triggering the devices.
Wild boar can thrive in captivity, though piglets grow slowly and poorly without their mothers. Products derived from wild boar include meat, hide and bristles. Apicius devotes a whole chapter to the cooking of boar meat, providing 10 recipes involving roasting, boiling and what sauces to use. The Romans usually served boar meat with garum. Boar's head was the centrepiece of most medieval Christmas celebrations among the nobility. Although growing in popularity as a captive-bred source of food, the wild boar takes longer to mature than most domestic pigs and it is usually smaller and produces less meat. Nevertheless, wild boar meat is leaner and healthier than pork, being of higher nutritional value and having a much higher concentration of essential amino acids. Most meat-dressing organizations agree that a boar carcass should yield of meat on average. Large specimens can yield of fat, with some giants yielding or more. A boar hide can measure and can yield of bristle and of underwool.
Crop and garbage raiding
Boars can be damaging to agriculture in situations where their natural habitat is sparse. Populations living on the outskirts of towns or farms can dig up potatoes and damage melons, watermelons and maize. However, they generally only encroach upon farms when natural food is scarce. In the Belovezh forest for example, 34–47% of the local boar population will enter fields in years of moderate availability of natural foods. While the role of boars in damaging crops is often exaggerated, cases are known of boar depredations causing famines, as was the case in Hachinohe, Japan in 1749, where 3,000 people died of what became known as the "wild boar famine". Still, within Japanese culture, the boar's status as vermin is expressed through its title as "king of pests" and the popular saying (addressed to young men in rural areas) "When you get married, choose a place with no wild boar."
In Central Europe, farmers typically repel boars through distraction or fright, while in Kazakhstan it is usual to employ guard dogs in plantations. However, research shows that when compared with other mitigation tactics, hunting is the only strategy to significantly reduce crop damage by boars. Although large boar populations can play an important role in limiting forest growth, they are also useful in keeping pest populations such as June bugs under control. The growth of urban areas and the corresponding decline in natural boar habitats has led to some sounders entering human habitations in search of food. As in natural conditions, sounders in peri-urban areas are matriarchal, though males tend to be much less represented and adults of both sexes can be up to 35% heavier than their forest-dwelling counterparts. As of 2010, at least 44 cities in 15 countries have experienced problems of some kind relating to the presence of habituated wild boar.
A 2023 study found that allowing wild pigs to forage on edible garbage in large regional landfills results in those animals getting physically large/heavier, having larger litters of piglets, and causing more wild pig-vehicle collisions in the vicinity of the landfill. The effects of letting these pigs scavenge in these landfills can present unique challenges to population management, control, public safety, and disease transmission. Wild pigs foraging on edible food waste in landfills has also been identified as a vector that facilitates the spread of African swine fever virus.
Attacks on humans
Actual attacks on humans are rare, but can be serious, resulting in penetrating injuries to the lower part of the body. They generally occur during the boars' rutting season from November to January, in agricultural areas bordering forests or on paths leading through forests. The animal typically attacks by charging and pointing its tusks towards the intended victim, with most injuries occurring on the thigh region. Once the initial attack is over, the boar steps back, takes position and attacks again if the victim is still moving, only ending once the victim is completely incapacitated.
Boar attacks on humans have been documented throughout history. The Romans and Ancient Greeks wrote of these attacks (Odysseus was wounded by a boar and Adonis was killed by one). A 2012 study compiling recorded attacks from 1825 to 2012 found accounts of 665 human victims of both wild boars and feral pigs, with the majority (19%) of attacks in the animal's native range occurring in India. Most of the attacks occurred in rural areas during the winter months in non-hunting contexts and were committed by solitary males.
Management
Managing wild boar is a pressing task in both native and invasive contexts as they can be disruptive to other systems when not addressed. Wild boar find their success through adaptation of daily patterns to circumvent threats. They avoid human contact through nocturnal lifestyles, despite the fact that they are not evolutionarily predisposed, and alter their diets substantially based on what is available. These "adaptive generalists", can survive in a variety of landscapes, making the prediction of their movement patterns and any potential close contact areas crucial to limiting damage. All of these qualities make them equally difficult to manage or limit.
Within Central Europe, the native habitat of the wild boar, there has been a push to re-evaluate interactions between wild boar and humans, with the priority of fostering positive engagement. Negative media and public perception of wild boar as "crop raiders" have made those living alongside them less willing to accept the economic damages of their behaviors, as wild boar are seen as pests. This media tone impacts management policy, with every 10 negative articles increasing wild boar policy activity by 6.7%. Contrary to this portrayal, wild boar, when managed well within their natural environments, can be a crucial part of forest ecosystems.
Defining the limits of proper management is difficult, but the exclusion of wild boar from rare environments is generally agreed upon, as when not properly managed, they can damage agricultural ventures and harm vulnerable plant life. These damages are estimated at $800 million yearly in environmental and financial costs for the United States alone. The breadth of this damage is due to prior inattention and lack of management tactics for extended lengths of time. Managing wild boar is a complex task, as it involves coordinating a combination of crop harvest techniques, fencing, toxic bait, corrals, and hunting. The most common tactic employed by private land owners in the United States is recreational hunting, however, this is generally not as effective on its own. Management strategies are most successful when they take into account reproduction, dispersion, and the differences between ideal resources for males and females.
According to a study, wild boars are causing soil disturbance that, among other problems, globally results in annual carbon dioxide emissions equivalent to that of ~1.1 million passenger vehicles (4.9 Mt, 0.01% of all GHG emissions as of 2022), implying that as of 2021 hunted boar meat – unlike other meat products – has beneficial effects on the environment even though the effect would diminish if boars are introduced for meat production and consistently retaining small populations of boars may be preferable.
| Biology and health sciences | Artiodactyla | null |
52316 | https://en.wikipedia.org/wiki/Mood%20disorder | Mood disorder | A mood disorder, also known as an affective disorder, is any of a group of conditions of mental and behavioral disorder where the main underlying characteristic is a disturbance in the person's mood. The classification is in the Diagnostic and Statistical Manual of Mental Disorders (DSM) and International Classification of Diseases (ICD).
Mood disorders fall into seven groups, including; abnormally elevated mood, such as mania or hypomania; depressed mood, of which the best-known and most researched is major depressive disorder (MDD) (alternatively known as clinical depression, unipolar depression, or major depression); and moods which cycle between mania and depression, known as bipolar disorder (BD) (formerly known as manic depression). There are several subtypes of depressive disorders or psychiatric syndromes featuring less severe symptoms such as dysthymic disorder (similar to MDD, but longer lasting and more persistent, though often milder) and cyclothymic disorder (similar to but milder than BD).
In some cases, more than one mood disorder can be present in an individual, like bipolar disorder and depressive disorder. Mood disorders may also be substance induced, or occur in response to a medical condition.
English psychiatrist Henry Maudsley proposed an overarching category of affective disorder. The term was then replaced by mood disorder, as the latter refers to the underlying or longitudinal emotional state, whereas the former refers to the external expression observed by others.
Classification
Depressive disorders
Major depressive disorder (MDD), commonly called major depression, unipolar depression, or clinical depression, wherein a person has one or more major depressive episodes. After a single episode, Major Depressive Disorder (single episode) would be diagnosed. After more than one episode, the diagnosis becomes Major Depressive Disorder (Recurrent). Depression without periods of mania is sometimes referred to as unipolar depression because the mood remains at the bottom "pole" and does not climb to the higher, manic "pole" as in bipolar disorder.
Individuals with a major depressive episode or major depressive disorder are at increased risk for suicide. Seeking help and treatment from a health professional dramatically reduces the individual's risk for suicide. Studies have demonstrated that asking if a depressed friend or family member has thought of committing suicide is an effective way of identifying those at risk, and it does not "plant" the idea or increase an individual's risk for suicide in any way. Epidemiological studies carried out in Europe suggest that, at this moment, roughly 8.5 percent of the world's population have a depressive disorder. No age group seems to be exempt from depression, and studies have found that depression appears in infants as young as 6 months old who have been separated from their mothers. However, there may be differences between cultures in prevalence of MDD due to cultural influences that "challenge the definition and diagnosis of psychiatric disorders", as seen in a study by Parker et al. that researched MDD in Chinese individuals. Depressive disorder (also known as depression) is a common mental disorder. It involves a depressed mood or loss of pleasure or interest in activities for long periods of time. Depression is different from regular mood changes and feelings about everyday life. It can affect all aspects of life. Depression can happen to anyone. People who have lived through abuse, severe losses, or other stressful events are more likely to develop depression.
Depressive disorder is frequent in primary care and general hospital practice but is often undetected. Unrecognized depressive disorder may slow recovery and worsen prognosis in physical illness, therefore it is important that all doctors be able to recognize the condition, treat the less severe cases, and identify those requiring specialist care.
Diagnosticians recognize several subtypes or course specifiers:
Atypical depression (AD) is characterized by mood reactivity (paradoxical anhedonia) and positivity, significant weight gain or increased appetite ("comfort eating"), excessive sleep or somnolence (hypersomnia), a sensation of heaviness in limbs known as leaden paralysis, and significant social impairment as a consequence of hypersensitivity to perceived interpersonal rejection. Difficulties in measuring this subtype have led to questions of its validity and prevalence.
Melancholic depression is characterized by a loss of pleasure (anhedonia) in most or all activities, a failure of reactivity to pleasurable stimuli, a quality of depressed mood more pronounced than that of grief or loss, a worsening of symptoms in the morning hours, early-morning waking, psychomotor retardation, excessive weight loss (not to be confused with anorexia nervosa), or excessive guilt.
Psychotic major depression (PMD), or simply psychotic depression, is the term for a major depressive episode, in particular of melancholic nature, wherein the patient experiences psychotic symptoms such as delusions or, less commonly, hallucinations. These are most commonly mood-congruent (content coincident with depressive themes).
Catatonic depression is a rare and severe form of major depression involving disturbances of motor behavior and other symptoms. Here, the person is mute and almost stuporous, and either is immobile or exhibits purposeless or even bizarre movements. Catatonic symptoms can also occur in schizophrenia or a manic episode, or can be due to neuroleptic malignant syndrome.
Postpartum depression (PPD) is listed as a course specifier in DSM-IV-TR; it refers to the intense, sustained and sometimes disabling depression experienced by women after giving birth. Postpartum depression, which affects 10–15% of women, typically sets in within three months of labor, and lasts as long as three months. It is quite common for women to experience a short-term feeling of tiredness and sadness in the first few weeks after giving birth; however, postpartum depression is different because it can cause significant hardship and impaired functioning at home, work, or school as well as, possibly, difficulty in relationships with family members, spouses, or friends, or even problems bonding with the newborn. In the treatment of postpartum major depressive disorders and other unipolar depressions in women who are breastfeeding, nortriptyline, paroxetine (Paxil), and sertraline (Zoloft) are in general considered to be the preferred medications. Women with personal or family histories of mood disorders are at particularly high risk of developing postpartum depression.
Premenstrual dysphoric disorder (PMDD) is a severe and disabling form of premenstrual syndrome affecting 3–8% of menstruating women. The disorder consists of a "cluster of affective, behavioral and somatic symptoms" that recur monthly during the luteal phase of the menstrual cycle. PMDD was added to the list of depressive disorders in the Diagnostic and Statistical Manual of Mental Disorders in 2013. The exact pathogenesis of the disorder is still unclear and is an active research topic. Treatment of PMDD relies largely on antidepressants that modulate serotonin levels in the brain via serotonin reuptake inhibitors as well as ovulation suppression using contraception.
Seasonal affective disorder (SAD), also known as "winter depression" or "winter blues", is a specifier. Some people have a seasonal pattern, with depressive episodes coming on in the autumn or winter, and resolving in spring. The diagnosis is made if at least two episodes have occurred in colder months with none at other times over a two-year period or longer. It is commonly hypothesised that people who live at higher latitudes tend to have less sunlight exposure in the winter and therefore experience higher rates of SAD, but the epidemiological support for this proposition is not strong (and latitude is not the only determinant of the amount of sunlight reaching the eyes in winter). It is said that this disorder can be treated by light therapy. SAD is also more prevalent in people who are younger and typically affects more females than males.
Dysthymia is a condition related to unipolar depression, where the same physical and cognitive problems are evident, but they are not as severe and tend to last longer (usually at least 2 years). The treatment of dysthymia is largely the same as for major depression, including antidepressant medications and psychotherapy.
Double depression can be defined as a fairly depressed mood (dysthymia) that lasts for at least two years and is punctuated by periods of major depression.
Unspecified Depressive Disorder is designated by the code 311 for depressive disorders. In the DSM-5, Unspecified Depressive Disorder encompasses symptoms that are characteristic of depressive disorders and cause significant impairment in functioning, but do not meet the criteria for the diagnosis of any specified depressive disorders. In the DSM-IV, this was called Depressive Disorder Not Otherwise Specified.
Depressive personality disorder (DPD) is a controversial psychiatric diagnosis that denotes a personality disorder with depressive features. Originally included in the DSM-II, depressive personality disorder was removed from the DSM-III and DSM-III-R. Recently, it has been reconsidered for reinstatement as a diagnosis. Depressive personality disorder is currently described in Appendix B in the DSM-IV-TR as worthy of further study.
Recurrent brief depression (RBD), distinguished from major depressive disorder primarily by differences in duration. People with RBD have depressive episodes about once per month, with individual episodes lasting less than two weeks and typically less than 2–3 days. Diagnosis of RBD requires that the episodes occur over the span of at least one year and, in female patients, independently of the menstrual cycle. People with clinical depression can develop RBD, and vice versa and both illnesses have similar risks.
Minor depressive disorder, or simply minor depression, which refers to a depression that does not meet full criteria for major depression but in which at least two symptoms are present for two weeks.
Bipolar disorders
Bipolar disorder (BD) (also called "manic depression" or "manic-depressive disorder"), an unstable emotional condition characterized by cycles of abnormal, persistent high mood (mania) and low mood (depression), which was formerly known as "manic depression" (and in some cases rapid cycling, mixed states, and psychotic symptoms). Subtypes include:
Bipolar I is distinguished by the presence or history of one or more manic episodes or mixed episodes with or without major depressive episodes. A depressive episode is not required for the diagnosis of Bipolar I Disorder, but depressive episodes are usually part of the course of the illness.
Bipolar II consisting of recurrent intermittent hypomanic and depressive episodes or mixed episodes.
Cyclothymia is a form of bipolar disorder, consisting of recurrent hypomanic and dysthymic episodes, but no full manic episodes or full major depressive episodes.
Bipolar disorder not otherwise specified (BD-NOS), sometimes called "sub-threshold" bipolar, indicates that the patient has some symptoms in the bipolar spectrum (e.g., manic and depressive symptoms) but does not fully qualify for any of the three formal bipolar DSM-IV diagnoses mentioned above.
It is estimated that roughly 1% of the adult population has bipolar I, a further 1% has bipolar II or cyclothymia, and somewhere between 2% and 5% percent have "sub-threshold" forms of bipolar disorder. Furthermore, the possibility of getting bipolar disorder when one parent is diagnosed with it is 15–30%. Risk, when both parents have it, is 50–75%. Also, while with bipolar siblings the risk is 15–25%, with identical twins it is about 70%.
Substance-induced
A mood disorder can be classified as substance-induced if its etiology can be traced to the direct physiologic effects of a psychoactive drug or other chemical substance, or if the development of the mood disorder occurred contemporaneously with substance intoxication or withdrawal. Also, an individual may have a mood disorder coexisting with a substance abuse disorder. Substance-induced mood disorders can have features of a manic, hypomanic, mixed, or depressive episode. Most substances can induce a variety of mood disorders. For example, stimulants such as amphetamine, methamphetamine, and cocaine can cause manic, hypomanic, mixed, and depressive episodes.
Alcohol-induced
High rates of major depressive disorder occur in heavy drinkers and those with alcoholism. Controversy has previously surrounded whether those who abused alcohol and developed depression were self-medicating their pre-existing depression. Recent research has concluded that, while this may be true in some cases, alcohol misuse directly causes the development of depression in a significant number of heavy drinkers. Participants studied were also assessed during stressful events in their lives and measured on a Feeling Bad Scale. Likewise, they were also assessed on their affiliation with deviant peers, unemployment, and their partner's substance use and criminal offending. High rates of suicide also occur in those who have alcohol-related problems. It is usually possible to differentiate between alcohol-related depression and depression that is not related to alcohol intake by taking a careful history of the patient. Depression and other mental health problems associated with alcohol misuse may be due to distortion of brain chemistry, as they tend to improve on their own after a period of abstinence.
Benzodiazepine-induced
Benzodiazepines, such as alprazolam, clonazepam, lorazepam and diazepam, can cause both depression and mania.
Benzodiazepines are a class of medication commonly used to treat anxiety, panic attacks and insomnia, and are also commonly misused and abused. Those with anxiety, panic and sleep problems commonly have negative emotions and thoughts, depression, suicidal ideations, and often have comorbid depressive disorders. While the anxiolytic and hypnotic effects of benzodiazepines may disappear as tolerance develops, depression and impulsivity with high suicidal risk commonly persist. These symptoms are "often interpreted as an exacerbation or as a natural evolution of previous disorders and the chronic use of sedatives is overlooked". Benzodiazepines do not prevent the development of depression, can exacerbate preexisting depression, can cause depression in those with no history of it, and can lead to suicide attempts. Risk factors for suicide and suicide attempts while using benzodiazepines include high dose prescriptions (even in those not misusing the medications), benzodiazepine intoxication, and underlying depression.
The long-term use of benzodiazepines may have a similar effect on the brain as alcohol, and are also implicated in depression. As with alcohol, the effects of benzodiazepine on neurochemistry, such as decreased levels of serotonin and norepinephrine, are believed to be responsible for the increased depression. Additionally, benzodiazepines can indirectly worsen mood by worsening sleep (i.e., benzodiazepine-induced sleep disorder). Like alcohol, benzodiazepines can put people to sleep but, while asleep, they disrupt sleep architecture: decreasing sleep time, delaying time to REM sleep, and decreasing deep sleep (the most restorative part of sleep for both energy and mood). Just as some antidepressants can cause or worsen anxiety in some patients due to being activating, benzodiazepines can cause or worsen depression due to being a central nervous system depressant—worsening thinking, concentration and problem solving (i.e., benzodiazepine-induced neurocognitive disorder). However, unlike antidepressants, in which the activating effects usually improve with continued treatment, benzodiazepine-induced depression is unlikely to improve until after stopping the medication.
In a long-term follow-up study of patients dependent on benzodiazepines, it was found that 10 people (20%) had taken drug overdoses while on chronic benzodiazepine medication despite only two people ever having had any pre-existing depressive disorder. A year after a gradual withdrawal program, no patients had taken any further overdoses.
Just as with intoxication and chronic use, benzodiazepine withdrawal can also cause depression. While benzodiazepine-induced depressive disorder may be exacerbated immediately after discontinuation of benzodiazepines, evidence suggests that mood significantly improves after the acute withdrawal period to levels better than during use. Depression resulting from withdrawal from benzodiazepines usually subsides after a few months but in some cases may persist for 6–12 months.
Due to another medical condition
"Mood disorder due to a general medical condition" is used to describe manic or depressive episodes which occur secondary to a medical condition. There are many medical conditions that can trigger mood episodes, including neurological disorders (e.g. dementias), hearing loss and associated disorders (e.g. tinnitus or hyperacusis), metabolic disorders (e.g. electrolyte disturbances), gastrointestinal diseases (e.g. cirrhosis), endocrine disease (e.g. thyroid abnormalities), cardiovascular disease (e.g. heart attack), pulmonary disease (e.g. chronic obstructive pulmonary disease), cancer, autoimmune diseases (e.g. multiple sclerosis), and pregnancy.
Not otherwise specified
Mood disorder not otherwise specified (MD-NOS) is a mood disorder that is impairing but does not fit in with any of the other officially specified diagnoses. In the DSM-IV MD-NOS is described as "any mood disorder that does not meet the criteria for a specific disorder." MD-NOS is not used as a clinical description but as a statistical concept for filing purposes. The diagnosis of MD-NOS does not exist in the DSM-5, however the diagnoses of unspecified depressive disorder and unspecified bipolar disorder are in the DSM-5.
Most cases of MD-NOS represent hybrids between mood and anxiety disorders, such as mixed anxiety-depressive disorder or atypical depression. An example of an instance of MD-NOS is being in minor depression frequently during various intervals, such as once every month or once in three days. There is a risk for MD-NOS not to get noticed, and for that reason not to get treated.
Causes
Meta-analyses show that high scores on the personality domain neuroticism are a strong predictor for the development of mood disorders. A number of authors have also suggested that mood disorders are an evolutionary adaptation (see also evolutionary psychiatry). A low or depressed mood can increase an individual's ability to cope with situations in which the effort to pursue a major goal could result in danger, loss, or wasted effort. In such situations, low motivation may give an advantage by inhibiting certain actions. This theory helps to explain why negative life incidents precede depression in around 80 percent of cases, and why they so often strike people during their peak reproductive years. These characteristics would be difficult to understand if depression were a dysfunction.
A depressed mood is a predictable response to certain types of life occurrences, such as loss of status, divorce, or death of a child or spouse. These are events that signal a loss of reproductive ability or potential, or that did so in humans' ancestral environment. A depressed mood can be seen as an adaptive response, in the sense that it causes an individual to turn away from the earlier (and reproductively unsuccessful) modes of behavior.
A depressed mood is common during illnesses, such as influenza. It has been argued that this is an evolved mechanism that assists the individual in recovering by limiting their physical activity. The occurrence of low-level depression during the winter months, or seasonal affective disorder, may have been adaptive in the past, by limiting physical activity at times when food was scarce. It is argued that humans have retained the instinct to experience low mood during the winter months, even if the availability of food is no longer determined by the weather.
Much of what is known about the genetic influence of clinical depression is based upon research that has been done with identical twins. Identical twins have exactly the same genetic code. It has been found that when one identical twin becomes depressed the other will also develop clinical depression approximately 76% of the time. When identical twins are raised apart from each other, they will both become depressed about 67% of the time. Because both twins become depressed at such a high rate, the implication is that there is a strong genetic influence. If it happened that when one twin becomes clinically depressed the other always develops depression, then clinical depression would likely be entirely genetic.
Bipolar disorder is also considered a mood disorder and it is hypothesized that it might be caused by mitochondrial dysfunction.
Sex differences
Mood disorders, specifically stress-related mood disorders such as anxiety and depression, have been shown to have differing rates of diagnosis based on sex. In the United States, women are two times more likely than men to be diagnosed with a stress-related mood disorder. Underlying these sex differences, studies have shown a dysregulation of stress-responsive neuroendocrine function causing an increase in the likelihood of developing these affective disorders. Overactivation of the hypothalamic-pituitary-adrenal (HPA) axis could provide potential insight into how these sex differences arise. Neuropeptide corticotropin-releasing factor (CRF) is released from the paraventricular nucleus (PVN) of the hypothalamus, stimulating adrenocorticotropic hormone (ACTH) release into the blood stream. From here ACTH triggers the release of glucocorticoids such as cortisol from the adrenal cortex. Cortisol, known as the main stress hormone, creates a negative feedback loop back to the hypothalamus to deactivate the stress response. When a constant stressor is present, the HPA axis remains overactivated and cortisol is constantly produced. This chronic stress is associated with sustained CRF release, resulting in the increased production of anxiety- and depressive-like behaviors and serving as a potential mechanism for differences in prevalence between men and women.
Diagnosis
DSM-5
The DSM-5, released in May 2013, separates the mood disorder chapter from the DSM-IV-TR into two sections: Depressive and related disorders and bipolar and related disorders. Bipolar disorders fall in between depressive disorders and schizophrenia spectrum and related disorders "in recognition of their place as a bridge between the two diagnostic classes in terms of symptomatology, family history and genetics" (Ref. 1, p 123). Bipolar disorders underwent a few changes in the DSM-5, most notably the addition of more specific symptomology related to hypomanic and mixed manic states. Depressive disorders underwent the most changes, the addition of three new disorders: disruptive mood dysregulation disorder, persistent depressive disorder (previously dysthymia), and premenstrual dysphoric disorder (previously in appendix B, the section for disorders needing further research). Disruptive mood dysregulation disorder is meant as a diagnosis for children and adolescents who would normally be diagnosed with bipolar disorder as a way to limit the bipolar diagnosis in this age cohort. Major depressive disorder (MDD) also underwent a notable change, in that the bereavement clause has been removed. Those previously exempt from a diagnosis of MDD due to bereavement are now candidates for the MDD diagnosis.
Treatment
There are different types of treatments available for mood disorders, such as therapy and medications. Behaviour therapy, cognitive behaviour therapy and interpersonal therapy have all shown to be potentially beneficial in depression. Major depressive disorder medications usually include antidepressants; a combination of antidepressants and cognitive behavioral therapy has shown to be more effective than one treatment alone. Bipolar disorder medications can consist of antipsychotics, mood stabilizers, anticonvulsants and/or lithium. Lithium specifically has been proven to reduce suicide and all causes of mortality in people with mood disorders. If mitochondrial dysfunction or mitochondrial diseases are the cause of mood disorders like bipolar disorder, then it has been hypothesized that N-acetyl-cysteine (NAC), acetyl-L-carnitine (ALCAR), S-adenosylmethionine (SAMe), coenzyme Q10 (CoQ10), alpha-lipoic acid (ALA), creatine monohydrate (CM), and melatonin could be potential treatment options. In determining treatment, there are many types of depression scales that are used. One of the depression scales is a self-report scale called Beck Depression Inventory (BDI). Another scale is the Hamilton Depression Rating Scale (HAMD). HAMD is a clinical rating scale in which the patient is rated based on clinician observation. The Center for Epidemiologic Studies Depression Scale (CES-D) is a scale for depression symptoms that applies to the general population. This scale is typically used in research and not for self-reports. The PHQ-9 which stands for Patient-Health Questionnaire-9 questions, is a self-report as well. Finally, the Mood Disorder Questionnaire (MDQ) evaluates bipolar disorder.
Epidemiology
According to a substantial number of epidemiology studies conducted, women are twice as likely to develop certain mood disorders, such as major depression. Although there is an equal number of men and women diagnosed with bipolar II disorder, women have a slightly higher frequency of the disorder.
In 2011, mood disorders were the most common reason for hospitalization among children aged 1–17 years in the United States, with approximately 112,000 stays. Mood disorders were top principal diagnosis for Medicaid super-utilizers in the United States in 2012. Further, a study of 18 states found that mood disorders accounted for the highest number of hospital readmissions among Medicaid patients and the uninsured, with 41,600 Medicaid patients and 12,200 uninsured patients being readmitted within 30 days of their index stay—a readmission rate of 19.8 per 100 admissions and 12.7 per 100 admissions, respectively. In 2012, mood and other behavioral health disorders were the most common diagnoses for Medicaid-covered and uninsured hospital stays in the United States (6.1% of Medicaid stays and 5.2% of uninsured stays).
A study conducted in 1988 to 1994 amongst young American adults involved a selection of demographic and health characteristics. A population-based sample of 8,602 men and women ages 17–39 years participated. Lifetime prevalence were estimated based on six mood measures:
major depressive episode (MDE) 8.6%,
major depressive disorder with severity (MDE-s) 7.7%,
dysthymia 6.2%,
MDE-s with dysthymia 3.4%,
any bipolar disorder 1.6%, and
any mood disorder 11.5%.
Research
Kay Redfield Jamison and others have explored the possible links between mood disorders – especially bipolar disorder – and creativity. It has been proposed that a "ruminating personality type may contribute to both [mood disorders] and art."
Jane Collingwood notes an Oregon State University study that:looked at the occupational status of a large group of typical patients and found that 'those with bipolar illness appear to be disproportionately concentrated in the most creative occupational category.' They also found that the likelihood of 'engaging in creative activities on the job' is significantly higher for bipolar than nonbipolar workers.In Liz Paterek's article "Bipolar Disorder and the Creative Mind" she wrote:Memory and creativity are related to mania. Clinical studies have shown that those in a manic state will rhyme, find synonyms, and use alliteration more than controls. This mental fluidity could contribute to an increase in creativity. Moreover, mania creates increases in productivity and energy. Those in a manic state are more emotionally sensitive and show less inhibition about attitudes, which could create greater expression. Studies performed at Harvard looked into the amount of original thinking in solving creative tasks. Bipolar individuals, whose disorder was not severe, tended to show greater degrees of creativity.The relationship between depression and creativity appears to be especially strong among poets.
| Biology and health sciences | Mental disorders | Health |
52358 | https://en.wikipedia.org/wiki/Imaginary%20unit | Imaginary unit | The imaginary unit or unit imaginary number () is a mathematical constant that is a solution to the quadratic equation Although there is no real number with this property, can be used to extend the real numbers to what are called complex numbers, using addition and multiplication. A simple example of the use of in a complex number is
Imaginary numbers are an important mathematical concept; they extend the real number system to the complex number system in which at least one root for every nonconstant polynomial exists (see Algebraic closure and Fundamental theorem of algebra). Here, the term "imaginary" is used because there is no real number having a negative square.
There are two complex square roots of and , just as there are two complex square roots of every real number other than zero (which has one double square root).
In contexts in which use of the letter is ambiguous or problematic, the letter is sometimes used instead. For example, in electrical engineering and control systems engineering, the imaginary unit is normally denoted by instead of , because is commonly used to denote electric current.
Terminology
Square roots of negative numbers are called imaginary because in early-modern mathematics, only what are now called real numbers, obtainable by physical measurements or basic arithmetic, were considered to be numbers at all – even negative numbers were treated with skepticism – so the square root of a negative number was previously considered undefined or nonsensical. The name imaginary is generally credited to René Descartes, and Isaac Newton used the term as early as 1670. The notation was introduced by Leonhard Euler.
A unit is an undivided whole, and unity or the unit number is the number one ().
Definition
The imaginary unit is defined solely by the property that its square is −1:
With defined this way, it follows directly from algebra that and are both square roots of −1.
Although the construction is called "imaginary", and although the concept of an imaginary number may be intuitively more difficult to grasp than that of a real number, the construction is valid from a mathematical standpoint. Real number operations can be extended to imaginary and complex numbers, by treating as an unknown quantity while manipulating an expression (and using the definition to replace any occurrence of with ). Higher integral powers of are thus
and so on, cycling through the four values , , , and . As with any non-zero real number,
As a complex number, can be represented in rectangular form as , with a zero real component and a unit imaginary component. In polar form, can be represented as (or just ), with an absolute value (or magnitude) of 1 and an argument (or angle) of radians. (Adding any integer multiple of to this angle works as well.) In the complex plane, which is a special interpretation of a Cartesian plane, is the point located one unit from the origin along the imaginary axis (which is orthogonal to the real axis).
vs.
Being a quadratic polynomial with no multiple root, the defining equation has distinct solutions, which are equally valid and which happen to be additive and multiplicative inverses of each other. Although the two solutions are distinct numbers, their properties are indistinguishable; there is no property that one has that the other does not. One of these two solutions is labelled (or simply ) and the other is labelled , though it is inherently ambiguous which is which.
The only differences between and arise from this labelling. For example, by convention is said to have an argument of and is said to have an argument of related to the convention of labelling orientations in the Cartesian plane relative to the positive -axis with positive angles turning anticlockwise in the direction of the positive -axis. Also, despite the signs written with them, neither nor is inherently positive or negative in the sense that real numbers are.
A more formal expression of this indistinguishability of and is that, although the complex field is unique (as an extension of the real numbers) up to isomorphism, it is unique up to a isomorphism. That is, there are two field automorphisms of the complex numbers that keep each real number fixed, namely the identity and complex conjugation. For more on this general phenomenon, see Galois group.
Matrices
Using the concepts of matrices and matrix multiplication, complex numbers can be represented in linear algebra. The real unit and imaginary unit can be represented by any pair of matrices and satisfying and Then a complex number can be represented by the matrix and all of the ordinary rules of complex arithmetic can be derived from the rules of matrix arithmetic.
The most common choice is to represent and by the identity matrix and the matrix ,
Then an arbitrary complex number can be represented by:
More generally, any real-valued matrix with a trace of zero and a determinant of one squares to , so could be chosen for . Larger matrices could also be used; for example, could be represented by the identity matrix and could be represented by any of the Dirac matrices for spatial dimensions.
Root of
Polynomials (weighted sums of the powers of a variable) are a basic tool in algebra. Polynomials whose coefficients are real numbers form a ring, denoted an algebraic structure with addition and multiplication and sharing many properties with the ring of integers.
The polynomial has no real-number roots, but the set of all real-coefficient polynomials divisible by forms an ideal, and so there is a quotient ring This quotient ring is isomorphic to the complex numbers, and the variable expresses the imaginary unit.
Graphic representation
The complex numbers can be represented graphically by drawing the real number line as the horizontal axis and the imaginary numbers as the vertical axis of a Cartesian plane called the complex plane. In this representation, the numbers and are at the same distance from , with a right angle between them. Addition by a complex number corresponds to translation in the plane, while multiplication by a unit-magnitude complex number corresponds to rotation about the origin. Every similarity transformation of the plane can be represented by a complex-linear function
Geometric algebra
In the geometric algebra of the Euclidean plane, the geometric product or quotient of two arbitrary vectors is a sum of a scalar (real number) part and a bivector part. (A scalar is a quantity with no orientation, a vector is a quantity oriented like a line, and a bivector is a quantity oriented like a plane.) The square of any vector is a positive scalar, representing its length squared, while the square of any bivector is a negative scalar.
The quotient of a vector with itself is the scalar , and when multiplied by any vector leaves it unchanged (the identity transformation). The quotient of any two perpendicular vectors of the same magnitude, , which when multiplied rotates the divisor a quarter turn into the dividend, , is a unit bivector which squares to , and can thus be taken as a representative of the imaginary unit. Any sum of a scalar and bivector can be multiplied by a vector to scale and rotate it, and the algebra of such sums is isomorphic to the algebra of complex numbers. In this interpretation points, vectors, and sums of scalars and bivectors are all distinct types of geometric objects.
More generally, in the geometric algebra of any higher-dimensional Euclidean space, a unit bivector of any arbitrary planar orientation squares to , so can be taken to represent the imaginary unit .
Proper use
The imaginary unit was historically written and still is in some modern works. However, great care needs to be taken when manipulating formulas involving radicals. The radical sign notation is reserved either for the principal square root function, which is defined for only real or for the principal branch of the complex square root function. Attempting to apply the calculation rules of the principal (real) square root function to manipulate the principal branch of the complex square root function can produce false results:
Generally, the calculation rules
and
are guaranteed to be valid only for real, positive values of and .
When or is real but negative, these problems can be avoided by writing and manipulating expressions like , rather than . For a more thorough discussion, see the articles Square root and Branch point.
Properties
As a complex number, the imaginary unit follows all of the rules of complex arithmetic.
Imaginary integers and imaginary numbers
When the imaginary unit is repeatedly added or subtracted, the result is some integer times the imaginary unit, an imaginary integer; any such numbers can be added and the result is also an imaginary integer:
Thus, the imaginary unit is the generator of a group under addition, specifically an infinite cyclic group.
The imaginary unit can also be multiplied by any arbitrary real number to form an imaginary number. These numbers can be pictured on a number line, the imaginary axis, which as part of the complex plane is typically drawn with a vertical orientation, perpendicular to the real axis which is drawn horizontally.
Gaussian integers
Integer sums of the real unit and the imaginary unit form a square lattice in the complex plane called the Gaussian integers. The sum, difference, or product of Gaussian integers is also a Gaussian integer:
Quarter-turn rotation
When multiplied by the imaginary unit , any arbitrary complex number in the complex plane is rotated by a quarter turn or ) anticlockwise. When multiplied by , any arbitrary complex number is rotated by a quarter turn clockwise. In polar form:
In rectangular form,
Integer powers
The powers of repeat in a cycle expressible with the following pattern, where is any integer:
Thus, under multiplication, is a generator of a cyclic group of order 4, a discrete subgroup of the continuous circle group of the unit complex numbers under multiplication.
Written as a special case of Euler's formula for an integer ,
With a careful choice of branch cuts and principal values, this last equation can also apply to arbitrary complex values of , including cases like .
Roots
Just like all nonzero complex numbers, has two distinct square roots which are additive inverses. In polar form, they are
In rectangular form, they are
Squaring either expression yields
The three cube roots of are
For a general positive integer , the -th roots of are, for
The value associated with is the principal -th root of . The set of roots equals the corresponding set of roots of unity rotated by the principal -th root of . These are the vertices of a regular polygon inscribed within the complex unit circle.
Exponential and logarithm
The complex exponential function relates complex addition in the domain to complex multiplication in the codomain. Real values in the domain represent scaling in the codomain (multiplication by a real scalar) with representing multiplication by , while imaginary values in the domain represent rotation in the codomain (multiplication by a unit complex number) with representing a rotation by radian. The complex exponential is thus a periodic function in the imaginary direction, with period and image at points for all integers , a real multiple of the lattice of imaginary integers.
The complex exponential can be broken into even and odd components, the hyperbolic functions and or the trigonometric functions and :
Euler's formula decomposes the exponential of an imaginary number representing a rotation:
This fact can be used to demonstrate, among other things, the apparently counterintuitive result that is a real number.
The quotient with appropriate scaling, can be represented as an infinite partial fraction decomposition as the sum of reciprocal functions translated by imaginary integers:
Other functions based on the complex exponential are well-defined with imaginary inputs. For example, a number raised to the power is:
Because the exponential is periodic, its inverse the complex logarithm is a multi-valued function, with each complex number in the domain corresponding to multiple values in the codomain, separated from each-other by any integer multiple of One way of obtaining a single-valued function is to treat the codomain as a cylinder, with complex values separated by any integer multiple of treated as the same value; another is to take the domain to be a Riemann surface consisting of multiple copies of the complex plane stitched together along the negative real axis as a branch cut, with each branch in the domain corresponding to one infinite strip in the codomain. Functions depending on the complex logarithm therefore depend on careful choice of branch to define and evaluate clearly.
For example, if one chooses any branch where then when is a positive real number,
Factorial
The factorial of the imaginary unit is most often given in terms of the gamma function evaluated at :
The magnitude and argument of this number are:
| Mathematics | Basics | null |
52376 | https://en.wikipedia.org/wiki/Axiom%20of%20extensionality | Axiom of extensionality | The axiom of extensionality, also called the axiom of extent, is an axiom used in many forms of axiomatic set theory, such as Zermelo–Fraenkel set theory. The axiom defines what a set is. Informally, the axiom means that the two sets A and B are equal if and only if A and B have the same members.
In ZF set theory
In the formal language of the Zermelo–Fraenkel axioms, the axiom reads:
or in words:
If the sets and have the same members, then they are the same set.
In , all members of sets are themselves sets, but not in set theory with urelements. The axiom's usefulness can be seen from the fact that, if one accepts that , where is a set and is a formula that occurs free in but doesn't, then the axiom assures that there is a unique set whose members are precisely whatever objects (urelements or sets, as the case may be) satisfy the formula .
The converse of the axiom, , follows from the substitution property of equality. Despite this, the axiom is sometimes given directly as a biconditional, i.e., as .
In NF set theory
Quine's New Foundations (NF) set theory, in Quine's original presentations of it, treats the symbol for equality or identity as shorthand either for "if a set contains the left side of the equals sign as a member, then it also contains the right side of the equals sign as a member" (as defined in 1937), or for "an object is an element of the set on the left side of the equals sign if, and only if, it is also an element of the set on the right side of the equals sign" (as defined in 1951). That is, is treated as shorthand either for , as in the original 1937 paper, or for , as in Quine's Mathematical Logic (1951). The second version of the definition is exactly equivalent to the antecedent of the ZF axiom of extensionality, and the first version of the definition is still very similar to it. By contrast, however, the ZF set theory takes the symbol for identity or equality as a primitive symbol of the formal language, and defines the axiom of extensionality in terms of it. (In this paragraph, the statements of both versions of the definition were paraphrases, and quotation marks were only used to set the statements apart.)
In Quine's New Foundations for Mathematical Logic (1937), the original paper of NF, the name "principle of extensionality" is given to the postulate P1, , which, for readability, may be restated as . The definition D8, which defines the symbol for identity or equality, defines as shorthand for . In his Mathematical Logic (1951), having already developed quasi-quotation, Quine defines as shorthand for (definition D10), and does not define an axiom or principle "of extensionality" at all.
Thomas Forster, however, ignores these fine distinctions, and considers NF to accept the axiom of extensionality in its ZF form.
In ZU set theory
In the Scott–Potter (ZU) set theory, the "extensionality principle" is given as a theorem rather than an axiom, which is proved from the definition of a "collection".
In set theory with ur-elements
An ur-element is a member of a set that is not itself a set.
In the Zermelo–Fraenkel axioms, there are no ur-elements, but they are included in some alternative axiomatisations of set theory.
Ur-elements can be treated as a different logical type from sets; in this case, makes no sense if is an ur-element, so the axiom of extensionality simply applies only to sets.
Alternatively, in untyped logic, we can require to be false whenever is an ur-element.
In this case, the usual axiom of extensionality would then imply that every ur-element is equal to the empty set.
To avoid this consequence, we can modify the axiom of extensionality to apply only to nonempty sets, so that it reads:
That is:
Given any set A and any set B, if A is a nonempty set (that is, if there exists a member X of A), then if A and B have precisely the same members, then they are equal.
Yet another alternative in untyped logic is to define itself to be the only element of
whenever is an ur-element. While this approach can serve to preserve the axiom of extensionality, the axiom of regularity will need an adjustment instead.
| Mathematics | Axiomatic systems | null |
52385 | https://en.wikipedia.org/wiki/Axiom%20of%20pairing | Axiom of pairing | In axiomatic set theory and the branches of logic, mathematics, and computer science that use it, the axiom of pairing is one of the axioms of Zermelo–Fraenkel set theory. It was introduced by as a special case of his axiom of elementary sets.
Formal statement
In the formal language of the Zermelo–Fraenkel axioms, the axiom reads:
In words:
Given any object A and any object B, there is a set C such that, given any object D, D is a member of C if and only if D is equal to A or D is equal to B.
Consequences
As noted, what the axiom is saying is that, given two objects A and B, we can find a set C whose members are exactly A and B.
We can use the axiom of extensionality to show that this set C is unique. We call the set C the pair of A and B, and denote it {A,B}. Thus the essence of the axiom is:
Any two objects have a pair.
The set {A,A} is abbreviated {A}, called the singleton containing A. Note that a singleton is a special case of a pair. Being able to construct a singleton is necessary, for example, to show the non-existence of the infinitely descending chains from the Axiom of regularity.
The axiom of pairing also allows for the definition of ordered pairs. For any objects and , the ordered pair is defined by the following:
Note that this definition satisfies the condition
Ordered n-tuples can be defined recursively as follows:
Alternatives
Non-independence
The axiom of pairing is generally considered uncontroversial, and it or an equivalent appears in just about any axiomatization of set theory. Nevertheless, in the standard formulation of the Zermelo–Fraenkel set theory, the axiom of pairing follows from the axiom schema of replacement applied to any given set with two or more elements, and thus it is sometimes omitted. The existence of such a set with two elements, such as { {}, { {} } }, can be deduced either from the axiom of empty set and the axiom of power set or from the axiom of infinity.
In the absence of some of the stronger ZFC axioms, the axiom of pairing can still, without loss, be introduced in weaker forms.
Weaker
In the presence of standard forms of the axiom schema of separation we can replace the axiom of pairing by its weaker version:
.
This weak axiom of pairing implies that any given objects and are members of some set . Using the axiom schema of separation we can construct the set whose members are exactly and .
Another axiom which implies the axiom of pairing in the presence of the axiom of empty set is the axiom of adjunction
.
It differs from the standard one by use of instead of . Using {} for A and x for B, we get {x} for C. Then use {x} for A and y for B, getting {x,y} for C. One may continue in this fashion to build up any finite set. And this could be used to generate all hereditarily finite sets without using the axiom of union.
Stronger
Together with the axiom of empty set and the axiom of union, the axiom of
pairing can be generalised to the following schema:
that is:
Given any finite number of objects A1 through An, there is a set C whose members are precisely A1 through An.
This set C is again unique by the axiom of extensionality, and is denoted {A1,...,An}.
Of course, we can't refer to a finite number of objects rigorously without already having in our hands a (finite) set to which the objects in question belong. Thus, this is not a single statement but instead a schema, with a separate statement for each natural number n.
The case n = 1 is the axiom of pairing with A = A1 and B = A1.
The case n = 2 is the axiom of pairing with A = A1 and B = A2.
The cases n > 2 can be proved using the axiom of pairing and the axiom of union multiple times.
For example, to prove the case n = 3, use the axiom of pairing three times, to produce the pair {A1,A2}, the singleton {A3}, and then the pair {{A1,A2},{A3}}.
The axiom of union then produces the desired result, {A1,A2,A3}. We can extend this schema to include n=0 if we interpret that case as the axiom of empty set.
Thus, one may use this as an axiom schema in the place of the axioms of empty set and pairing. Normally, however, one uses the axioms of empty set and pairing separately, and then proves this as a theorem schema. Note that adopting this as an axiom schema will not replace the axiom of union, which is still needed for other situations.
| Mathematics | Axiomatic systems | null |
52386 | https://en.wikipedia.org/wiki/Axiom%20schema%20of%20specification | Axiom schema of specification | In many popular versions of axiomatic set theory, the axiom schema of specification, also known as the axiom schema of separation (Aussonderungsaxiom), subset axiom, axiom of class construction, or axiom schema of restricted comprehension is an axiom schema. Essentially, it says that any definable subclass of a set is a set.
Some mathematicians call it the axiom schema of comprehension, although others use that term for unrestricted comprehension''', discussed below.
Because restricting comprehension avoided Russell's paradox, several mathematicians including Zermelo, Fraenkel, and Gödel considered it the most important axiom of set theory.
Statement
One instance of the schema is included for each formula φ in the language of set theory with free variables among x, w1, ..., wn, A. So B does not occur free in φ. In the formal language of set theory, the axiom schema is:
or in words:
Given any set A, there is a set B (a subset of A) such that, given any set x, x is a member of B if and only if x is a member of A and φ holds for x.
Note that there is one axiom for every such predicate φ; thus, this is an axiom schema.
To understand this axiom schema, note that the set B must be a subset of A. Thus, what the axiom schema is really saying is that, given a set A and a predicate , we can find a subset B of A whose members are precisely the members of A that satisfy . By the axiom of extensionality this set is unique. We usually denote this set using set-builder notation as . Thus the essence of the axiom is:
Every subclass of a set that is defined by a predicate is itself a set.
The preceding form of separation was introduced in 1930 by Thoralf Skolem as a refinement of a previous, non-first-order<ref>F. R. Drake, Set Theory: An Introduction to Large Cardinals (1974), pp.12--13. ISBN 0 444 10535 2.</ref> form by Zermelo. The axiom schema of specification is characteristic of systems of axiomatic set theory related to the usual set theory ZFC, but does not usually appear in radically different systems of alternative set theory. For example, New Foundations and positive set theory use different restrictions of the axiom of comprehension of naive set theory. The Alternative Set Theory of Vopenka makes a specific point of allowing proper subclasses of sets, called semisets. Even in systems related to ZFC, this scheme is sometimes restricted to formulas with bounded quantifiers, as in Kripke–Platek set theory with urelements.
Relation to the axiom schema of replacement
The axiom schema of specification is implied by the axiom schema of replacement together with the axiom of empty set.
The axiom schema of replacement says that, if a function is definable by a formula , then for any set , there exists a set :
.
To derive the axiom schema of specification, let be a formula and a set, and define the function such that if is true and if is false, where such that is true. Then the set guaranteed by the axiom schema of replacement is precisely the set required in the axiom schema of specification. If does not exist, then in the axiom schema of specification is the empty set, whose existence (i.e., the axiom of empty set) is then needed.
For this reason, the axiom schema of specification is left out of some axiomatizations of ZF (Zermelo–Fraenkel set theory), although some authors, despite the redundancy, include both. Regardless, the axiom schema of specification is notable because it was in Zermelo's original 1908 list of axioms, before Fraenkel invented the axiom of replacement in 1922. Additionally, if one takes ZFC set theory (i.e., ZF with the axiom of choice), removes the axiom of replacement and the axiom of collection, but keeps the axiom schema of specification, one gets the weaker system of axioms called ZC (i.e., Zermelo's axioms, plus the axiom of choice).
Unrestricted comprehension
The axiom schema of unrestricted comprehension reads:
that is:
This set is again unique, and is usually denoted as
In an unsorted material set theory, the axiom or rule of full or unrestricted comprehension''' says that for any property P, there exists a set {x | P(x)} of all objects satisfying P.This axiom schema was tacitly used in the early days of naive set theory, before a strict axiomatization was adopted. However, it was later discovered to lead directly to Russell's paradox, by taking to be (i.e., the property that set is not a member of itself). Therefore, no useful axiomatization of set theory can use unrestricted comprehension. Passing from classical logic to intuitionistic logic does not help, as the proof of Russell's paradox is intuitionistically valid.
Accepting only the axiom schema of specification was the beginning of axiomatic set theory. Most of the other Zermelo–Fraenkel axioms (but not the axiom of extensionality, the axiom of regularity, or the axiom of choice) then became necessary to make up for some of what was lost by changing the axiom schema of comprehension to the axiom schema of specification – each of these axioms states that a certain set exists, and defines that set by giving a predicate for its members to satisfy, i.e. it is a special case of the axiom schema of comprehension.
It is also possible to prevent the schema from being inconsistent by restricting which formulae it can be applied to, such as only stratified formulae in New Foundations (see below) or only positive formulae (formulae with only conjunction, disjunction, quantification and atomic formulae) in positive set theory. Positive formulae, however, typically are unable to express certain things that most theories can; for instance, there is no complement or relative complement in positive set theory.
In NBG class theory
In von Neumann–Bernays–Gödel set theory, a distinction is made between sets and classes. A class is a set if and only if it belongs to some class . In this theory, there is a theorem schema that reads
that is,
provided that the quantifiers in the predicate are restricted to sets.
This theorem schema is itself a restricted form of comprehension, which avoids Russell's paradox because of the requirement that be a set. Then specification for sets themselves can be written as a single axiom
that is,
or even more simply
In this axiom, the predicate is replaced by the class , which can be quantified over. Another simpler axiom which achieves the same effect is
that is,
In higher-order settings
In a typed language where we can quantify over predicates, the axiom schema of specification becomes a simple axiom. This is much the same trick as was used in the NBG axioms of the previous section, where the predicate was replaced by a class that was then quantified over.
In second-order logic and higher-order logic with higher-order semantics, the axiom of specification is a logical validity and does not need to be explicitly included in a theory.
In Quine's New Foundations
In the New Foundations approach to set theory pioneered by W. V. O. Quine, the axiom of comprehension for a given predicate takes the unrestricted form, but the predicates that may be used in the schema are themselves restricted. The predicate ( is not in ) is forbidden, because the same symbol appears on both sides of the membership symbol (and so at different "relative types"); thus, Russell's paradox is avoided. However, by taking to be , which is allowed, we can form a set of all sets. For details, see stratification.
| Mathematics | Axiomatic systems | null |
52387 | https://en.wikipedia.org/wiki/Axiom%20schema%20of%20replacement | Axiom schema of replacement | In set theory, the axiom schema of replacement is a schema of axioms in Zermelo–Fraenkel set theory (ZF) that asserts that the image of any set under any definable mapping is also a set. It is necessary for the construction of certain infinite sets in ZF.
The axiom schema is motivated by the idea that whether a class is a set depends only on the cardinality of the class, not on the rank of its elements. Thus, if one class is "small enough" to be a set, and there is a surjection from that class to a second class, the axiom states that the second class is also a set. However, because ZFC only speaks of sets, not proper classes, the schema is stated only for definable surjections, which are identified with their defining formulas.
Statement
Suppose is a definable binary relation (which may be a proper class) such that for every set there is a unique set such that holds. There is a corresponding definable function , where if and only if . Consider the (possibly proper) class defined such that for every set , if and only if there is an with . is called the image of under , and denoted or (using set-builder notation) .
The axiom schema of replacement states that if is a definable class function, as above, and is any set, then the image is also a set. This can be seen as a principle of smallness: the axiom states that if is small enough to be a set, then is also small enough to be a set. It is implied by the stronger axiom of limitation of size.
Because it is impossible to quantify over definable functions in first-order logic, one instance of the schema is included for each formula in the language of set theory with free variables among ; but is not free in . In the formal language of set theory, the axiom schema is:
For the meaning of , see uniqueness quantification.
For clarity, in the case of no variables , this simplifies to:
So whenever specifies a unique -to- correspondence, akin to a function on , then all reached this way can be collected into a set , akin to .
Applications
The axiom schema of replacement is not necessary for the proofs of most theorems of ordinary mathematics. Indeed, Zermelo set theory (Z) already can interpret second-order arithmetic and much of type theory in finite types, which in turn are sufficient to formalize the bulk of mathematics. Although the axiom schema of replacement is a standard axiom in set theory today, it is often omitted from systems of type theory and foundation systems in topos theory.
At any rate, the axiom schema drastically increases the strength of ZF, both in terms of the theorems it can prove - for example the sets shown to exist - and also in terms of its proof-theoretic consistency strength, compared to Z. Some important examples follow:
Using the modern definition due to von Neumann, proving the existence of any limit ordinal greater than ω requires the replacement axiom. The ordinal number ω·2 = ω + ω is the first such ordinal. The axiom of infinity asserts the existence of an infinite set ω = {0, 1, 2, ...}. One may hope to define ω·2 as the union of the sequence {ω, ω + 1, ω + 2,...}. However, arbitrary such classes of ordinals need not be sets - for example, the class of all ordinals is not a set. Replacement now allows one to replace each finite number n in ω with the corresponding ω + n, and thus guarantees that this class is a set. As a clarification, note that one can easily construct a well-ordered set that is isomorphic to ω·2 without resorting to replacement – simply take the disjoint union of two copies of ω, with the second copy greater than the first – but that this is not an ordinal since it is not totally ordered by inclusion.
Larger ordinals rely on replacement less directly. For example, ω1, the first uncountable ordinal, can be constructed as follows – the set of countable well orders exists as a subset of by separation and powerset (a relation on A is a subset of , and so an element of the power set . A set of relations is thus a subset of ). Replace each well-ordered set with its ordinal. This is the set of countable ordinals ω1, which can itself be shown to be uncountable. The construction uses replacement twice; once to ensure an ordinal assignment for each well ordered set and again to replace well ordered sets by their ordinals. This is a special case of the result of Hartogs number, and the general case can be proved similarly.
In light of the above, the existence of an assignment of an ordinal to every well-ordered set requires replacement as well. Similarly the von Neumann cardinal assignment which assigns a cardinal number to each set requires replacement, as well as axiom of choice.
For sets of tuples recursively defined as and for large , the set has too high of a rank for its existence to be provable from set theory with just the axiom of power set, choice and without replacement.
Similarly, Harvey Friedman showed that at least some instances of replacement are required to show that Borel games are determined. The proven result is Donald A. Martin's Borel determinacy theorem. A later, more careful analysis by Martin of the result showed that it only requires replacement for functions with domain an arbitrary countable ordinal.
ZF with replacement proves the consistency of Z, as the set Vω·2 is a model of Z whose existence can be proved in ZF. The cardinal number is the first one which can be shown to exist in ZF but not in Z. For clarification, note that Gödel's second incompleteness theorem shows that each of these theories contains a sentence, "expressing" the theory's own consistency, that is unprovable in that theory, if that theory is consistent - this result is often loosely expressed as the claim that neither of these theories can prove its own consistency, if it is consistent.
Relation to other axiom schemas
Simplifications
Some simplifications may be made to the axiom schema of replacement to obtain different equivalent versions. Azriel Lévy showed that a version of replacement with parameters removed, i.e. the following schema, is equivalent to the original form. In particular the equivalence holds in the presence of the axioms of extensionality, pairing, union and powerset.
Collection
The axiom schema of collection is closely related to and frequently confused with the axiom schema of replacement.
Over the remainder of the ZF axioms, it is equivalent to the axiom schema of replacement. The axiom of collection is stronger than replacement in the absence of the power set axiom or its constructive counterpart of ZF and is used in the framework of IZF, which lacks the law of excluded middle, instead of Replacement which is weaker.
While replacement can be read to say that the image of a function is a set, collection speaks about images of relations and then merely says that some superclass of the relation's image is a set.
In other words, the resulting set has no minimality requirement, i.e. this variant also lacks the uniqueness requirement on .
That is, the relation defined by is not required to be a function—some may correspond to many 's in . In this case, the image set whose existence is asserted must contain at least one such for each in the original set, with no guarantee that it will contain only one.
Suppose that the free variables of are among ; but neither nor is free in . Then the axiom schema is:
The axiom schema is sometimes stated without prior restrictions (apart from not occurring free in ) on the predicate, :
In this case, there may be elements in that are not associated to any other sets by . However, the axiom schema as stated requires that, if an element of is associated with at least one set , then the image set will contain at least one such . The resulting axiom schema is also called the axiom schema of boundedness.
Separation
The axiom schema of separation, the other axiom schema in ZFC, is implied by the axiom schema of replacement and the axiom of empty set. Recall that the axiom schema of separation includes
for each formula in the language of set theory in which is not free, i.e. that does not mention .
The proof is as follows: Either contains some element validating , or it does not. In the latter case, taking the empty set for fulfills the relevant instance of the axiom schema of separation and one is done. Otherwise, choose such a fixed in that validates . Now define for use with replacement. Using function notation for this predicate , it acts as the identity wherever is true and as the constant function wherever is false. By case analysis, the possible values are unique for any , meaning indeed constitutes a class function. In turn, the image of under , i.e. the class , is granted to be a set by the axiom of replacement. This precisely validates the axiom of separation.
This result shows that it is possible to axiomatize ZFC with a single infinite axiom schema. Because at least one such infinite schema is required (ZFC is not finitely axiomatizable), this shows that the axiom schema of replacement can stand as the only infinite axiom schema in ZFC if desired. Because the axiom schema of separation is not independent, it is sometimes omitted from contemporary statements of the Zermelo-Fraenkel axioms.
Separation is still important, however, for use in fragments of ZFC, because of historical considerations, and for comparison with alternative axiomatizations of set theory. A formulation of set theory that does not include the axiom of replacement will likely include some form of the axiom of separation, to ensure that its models contain a sufficiently rich collection of sets. In the study of models of set theory, it is sometimes useful to consider models of ZFC without replacement, such as the models in von Neumann's hierarchy.
The proof given above assumes the law of excluded middle for the proposition that is inhabited by a set validating , and for any when stipulating that the relation is functional. The axiom of separation is explicitly included in constructive set theory, or a bounded variant thereof.
Reflection
Lévy's reflection principle for ZFC is equivalent to the axiom of replacement, assuming the axiom of infinity. Lévy's principle is as follows:
For any and any first-order formula , there exists an such that .
This is a schema that consists of countably many statements, one for each formula . Here, means with all quantifiers bounded to , i.e. but with every instance of and replaced with and respectively.
History
The axiom schema of replacement was not part of Ernst Zermelo's 1908 axiomatisation of set theory (Z). Some informal approximation to it existed in Cantor's unpublished works, and it appeared again informally in Mirimanoff (1917).
Its publication by Abraham Fraenkel in 1922 is what makes modern set theory Zermelo-Fraenkel set theory (ZFC). The axiom was independently discovered and announced by Thoralf Skolem later in the same year (and published in 1923). Zermelo himself incorporated Fraenkel's axiom in his revised system he published in 1930, which also included as a new axiom von Neumann's axiom of foundation. Although it is Skolem's first order version of the axiom list that we use today, he usually gets no credit since each individual axiom was developed earlier by either Zermelo or Fraenkel. The phrase “Zermelo-Fraenkel set theory” was first used in print by von Neumann in 1928.
Zermelo and Fraenkel had corresponded heavily in 1921; the axiom of replacement was a major topic of this exchange. Fraenkel initiated correspondence with Zermelo sometime in March 1921. However, his letters before the one dated 6 May 1921 are lost. Zermelo first admitted to a gap in his system in a reply to Fraenkel dated 9 May 1921. On 10 July 1921, Fraenkel completed and submitted for publication a paper (published in 1922) that described his axiom as allowing arbitrary replacements: "If M is a set and each element of M is replaced by [a set or an urelement] then M turns into a set again" (parenthetical completion and translation by Ebbinghaus). Fraenkel's 1922 publication thanked Zermelo for helpful arguments. Prior to this publication, Fraenkel publicly announced his new axiom at a meeting of the German Mathematical Society held in Jena on 22 September 1921. Zermelo was present at this meeting; in the discussion following Fraenkel's talk he accepted the axiom of replacement in general terms, but expressed reservations regarding its extent.
Thoralf Skolem made public his discovery of the gap in Zermelo's system (the same gap that Fraenkel had found) in a talk he gave on 6 July 1922 at the 5th Congress of Scandinavian Mathematicians, which was held in Helsinki; the proceedings of this congress were published in 1923. Skolem presented a resolution in terms of first-order definable replacements: "Let U be a definite proposition that holds for certain pairs (a, b) in the domain B; assume further, that for every a there exists at most one b such that U is true. Then, as a ranges over the elements of a set Ma, b ranges over all elements of a set Mb." In the same year, Fraenkel wrote a review of Skolem's paper, in which Fraenkel simply stated that Skolem's considerations correspond to his own.
Zermelo himself never accepted Skolem's formulation of the axiom schema of replacement. At one point he called Skolem's approach “set theory of the impoverished”. Zermelo envisaged a system that would allow for large cardinals. He also objected strongly to the philosophical implications of countable models of set theory, which followed from Skolem's first-order axiomatization. According to the biography of Zermelo by Heinz-Dieter Ebbinghaus, Zermelo's disapproval of Skolem's approach marked the end of Zermelo's influence on the developments of set theory and logic.
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52432 | https://en.wikipedia.org/wiki/Xanthine | Xanthine | Xanthine ( or , from Ancient Greek for its yellowish-white appearance; archaically xanthic acid; systematic name 3,7-dihydropurine-2,6-dione) is a purine base found in most human body tissues and fluids, as well as in other organisms. Several stimulants are derived from xanthine, including caffeine, theophylline, and theobromine.
Xanthine is a product on the pathway of purine degradation.
It is created from guanine by guanine deaminase.
It is created from hypoxanthine by xanthine oxidoreductase.
It is also created from xanthosine by purine nucleoside phosphorylase.
Xanthine is subsequently converted to uric acid by the action of the xanthine oxidase enzyme.
Use and production
Xanthine is used as a drug precursor for human and animal medications, and is produced as a pesticide ingredient.
Clinical significance
Derivatives of xanthine (known collectively as xanthines) are a group of alkaloids commonly used for their effects as mild stimulants and as bronchodilators, notably in the treatment of asthma or influenza symptoms. In contrast to other, more potent stimulants like sympathomimetic amines, xanthines mainly act to oppose the actions of adenosine, and increase alertness in the central nervous system.
Toxicity
Methylxanthines (methylated xanthines), which include caffeine, aminophylline, IBMX, paraxanthine, pentoxifylline, theobromine, theophylline, and 7-methylxanthine (heteroxanthine), among others, affect the airways, increase heart rate and force of contraction, and at high concentrations can cause cardiac arrhythmias. In high doses, they can lead to convulsions that are resistant to anticonvulsants. Methylxanthines induce gastric acid and pepsin secretions in the gastrointestinal tract. Methylxanthines are metabolized by cytochrome P450 in the liver.
If swallowed, inhaled, or exposed to the eyes in high amounts, xanthines can be harmful, and they may cause an allergic reaction if applied topically.
Pharmacology
In in vitro pharmacological studies, xanthines act as both competitive nonselective phosphodiesterase inhibitors and nonselective adenosine receptor antagonists. Phosphodiesterase inhibitors raise intracellular cAMP, activate PKA, inhibit TNF-α synthesis, and leukotriene and reduce inflammation and innate immunity. Adenosine receptor antagonists inhibit sleepiness-inducing adenosine.
However, different analogues show varying potency at the numerous subtypes, and a wide range of synthetic xanthines (some nonmethylated) have been developed searching for compounds with greater selectivity for phosphodiesterase enzyme or adenosine receptor subtypes.
Pathology
People with rare genetic disorders, specifically xanthinuria and Lesch–Nyhan syndrome, lack sufficient xanthine oxidase and cannot convert xanthine to uric acid.
Possible formation in absence of life
Studies reported in 2008, based on 12C/13C isotopic ratios of organic compounds found in the Murchison meteorite, suggested that xanthine and related chemicals, including the RNA component uracil, have been formed extraterrestrially. In August 2011, a report, based on NASA studies with meteorites found on Earth, was published suggesting xanthine and related organic molecules, including the DNA and RNA components adenine and guanine, were found in outer space.
| Physical sciences | Alkaloids | Chemistry |
52552 | https://en.wikipedia.org/wiki/Axiom%20of%20power%20set | Axiom of power set | In mathematics, the axiom of power set is one of the Zermelo–Fraenkel axioms of axiomatic set theory. It guarantees for every set the existence of a set , the power set of , consisting precisely of the subsets of . By the axiom of extensionality, the set is unique.
The axiom of power set appears in most axiomatizations of set theory. It is generally considered uncontroversial, although constructive set theory prefers a weaker version to resolve concerns about predicativity.
Formal statement
The subset relation is not a primitive notion in formal set theory and is not used in the formal language of the Zermelo–Fraenkel axioms. Rather, the subset relation is defined in terms of set membership, . Given this, in the formal language of the Zermelo–Fraenkel axioms, the axiom of power set reads:
where y is the power set of x, z is any element of y, w is any member of z.
In English, this says:
Given any set x, there is a set y such that, given any set z, this set z is a member of y if and only if every element of z is also an element of x.
Consequences
The power set axiom allows a simple definition of the Cartesian product of two sets and :
Notice that
and, for example, considering a model using the Kuratowski ordered pair,
and thus the Cartesian product is a set since
One may define the Cartesian product of any finite collection of sets recursively:
The existence of the Cartesian product can be proved without using the power set axiom, as in the case of the Kripke–Platek set theory.
Limitations
The power set axiom does not specify what subsets of a set exist, only that there is a set containing all those that do. Not all conceivable subsets are guaranteed to exist. In particular, the power set of an infinite set would contain only "constructible sets" if the universe is the constructible universe but in other models of ZF set theory could contain sets that are not constructible.
| Mathematics | Axiomatic systems | null |
52553 | https://en.wikipedia.org/wiki/Axiom%20of%20union | Axiom of union | In axiomatic set theory, the axiom of union is one of the axioms of Zermelo–Fraenkel set theory. This axiom was introduced by Ernst Zermelo.
Informally, the axiom states that for each set x there is a set y whose elements are precisely the elements of the elements of x.
Formal statement
In the formal language of the Zermelo–Fraenkel axioms, the axiom reads:
or in words:
Given any set A, there is a set B such that, for any element c, c is a member of B if and only if there is a set D such that c is a member of D and D is a member of A.
or, more simply:
For any set , there is a set which consists of just the elements of the elements of that set .
Relation to Pairing
The axiom of union allows one to unpack a set of sets and thus create a flatter set.
Together with the axiom of pairing, this implies that for any two sets, there is a set (called their union) that contains exactly the elements of the two sets.
Relation to Replacement
The axiom of replacement allows one to form many unions, such as the union of two sets.
However, in its full generality, the axiom of union is independent from the rest of the ZFC-axioms:
Replacement does not prove the existence of the union of a set of sets if the result contains an unbounded number of cardinalities.
Together with the axiom schema of replacement, the axiom of union implies that one can form the union of a family of sets indexed by a set.
Relation to Separation
In the context of set theories which include the axiom of separation, the axiom of union is sometimes stated in a weaker form which only produces a superset of the union of a set. For example, Kunen states the axiom as
which is equivalent to
Compared to the axiom stated at the top of this section, this variation asserts only one direction of the implication, rather than both directions.
Relation to Intersection
There is no corresponding axiom of intersection. If is a nonempty set containing , it is possible to form the intersection using the axiom schema of specification as
,
so no separate axiom of intersection is necessary. (If A is the empty set, then trying to form the intersection of A as
{c: for all D in A, c is in D}
is not permitted by the axioms. Moreover, if such a set existed, then it would contain every set in the "universe", but the notion of a universal set is antithetical to Zermelo–Fraenkel set theory.)
| Mathematics | Axiomatic systems | null |
52564 | https://en.wikipedia.org/wiki/Partial%20differential%20equation | Partial differential equation | In mathematics, a partial differential equation (PDE) is an equation which involves a multivariable function and one or more of its partial derivatives.
The function is often thought of as an "unknown" that solves the equation, similar to how is thought of as an unknown number solving, e.g., an algebraic equation like . However, it is usually impossible to write down explicit formulae for solutions of partial differential equations. There is correspondingly a vast amount of modern mathematical and scientific research on methods to numerically approximate solutions of certain partial differential equations using computers. Partial differential equations also occupy a large sector of pure mathematical research, in which the usual questions are, broadly speaking, on the identification of general qualitative features of solutions of various partial differential equations, such as existence, uniqueness, regularity and stability. Among the many open questions are the existence and smoothness of solutions to the Navier–Stokes equations, named as one of the Millennium Prize Problems in 2000.
Partial differential equations are ubiquitous in mathematically oriented scientific fields, such as physics and engineering. For instance, they are foundational in the modern scientific understanding of sound, heat, diffusion, electrostatics, electrodynamics, thermodynamics, fluid dynamics, elasticity, general relativity, and quantum mechanics (Schrödinger equation, Pauli equation etc.). They also arise from many purely mathematical considerations, such as differential geometry and the calculus of variations; among other notable applications, they are the fundamental tool in the proof of the Poincaré conjecture from geometric topology.
Partly due to this variety of sources, there is a wide spectrum of different types of partial differential equations, where the meaning of a solution depends on the context of the problem, and methods have been developed for dealing with many of the individual equations which arise. As such, it is usually acknowledged that there is no "universal theory" of partial differential equations, with specialist knowledge being somewhat divided between several essentially distinct subfields.
Ordinary differential equations can be viewed as a subclass of partial differential equations, corresponding to functions of a single variable. Stochastic partial differential equations and nonlocal equations are, as of 2020, particularly widely studied extensions of the "PDE" notion. More classical topics, on which there is still much active research, include elliptic and parabolic partial differential equations, fluid mechanics, Boltzmann equations, and dispersive partial differential equations.
Introduction
A function of three variables is "harmonic" or "a solution of the Laplace equation" if it satisfies the condition
Such functions were widely studied in the 19th century due to their relevance for classical mechanics, for example the equilibrium temperature distribution of a homogeneous solid is a harmonic function. If explicitly given a function, it is usually a matter of straightforward computation to check whether or not it is harmonic. For instance
and
are both harmonic while
is not. It may be surprising that the two examples of harmonic functions are of such strikingly different form. This is a reflection of the fact that they are not, in any immediate way, special cases of a "general solution formula" of the Laplace equation. This is in striking contrast to the case of ordinary differential equations (ODEs) roughly similar to the Laplace equation, with the aim of many introductory textbooks being to find algorithms leading to general solution formulas. For the Laplace equation, as for a large number of partial differential equations, such solution formulas fail to exist.
The nature of this failure can be seen more concretely in the case of the following PDE: for a function of two variables, consider the equation
It can be directly checked that any function of the form , for any single-variable functions and whatsoever, will satisfy this condition. This is far beyond the choices available in ODE solution formulas, which typically allow the free choice of some numbers. In the study of PDEs, one generally has the free choice of functions.
The nature of this choice varies from PDE to PDE. To understand it for any given equation, existence and uniqueness theorems are usually important organizational principles. In many introductory textbooks, the role of existence and uniqueness theorems for ODE can be somewhat opaque; the existence half is usually unnecessary, since one can directly check any proposed solution formula, while the uniqueness half is often only present in the background in order to ensure that a proposed solution formula is as general as possible. By contrast, for PDE, existence and uniqueness theorems are often the only means by which one can navigate through the plethora of different solutions at hand. For this reason, they are also fundamental when carrying out a purely numerical simulation, as one must have an understanding of what data is to be prescribed by the user and what is to be left to the computer to calculate.
To discuss such existence and uniqueness theorems, it is necessary to be precise about the domain of the "unknown function". Otherwise, speaking only in terms such as "a function of two variables", it is impossible to meaningfully formulate the results. That is, the domain of the unknown function must be regarded as part of the structure of the PDE itself.
The following provides two classic examples of such existence and uniqueness theorems. Even though the two PDE in question are so similar, there is a striking difference in behavior: for the first PDE, one has the free prescription of a single function, while for the second PDE, one has the free prescription of two functions.
Let denote the unit-radius disk around the origin in the plane. For any continuous function on the unit circle, there is exactly one function on such that and whose restriction to the unit circle is given by .
For any functions and on the real line , there is exactly one function on such that and with and for all values of .
Even more phenomena are possible. For instance, the following PDE, arising naturally in the field of differential geometry, illustrates an example where there is a simple and completely explicit solution formula, but with the free choice of only three numbers and not even one function.
If is a function on with then there are numbers , , and with .
In contrast to the earlier examples, this PDE is nonlinear, owing to the square roots and the squares. A linear PDE is one such that, if it is homogeneous, the sum of any two solutions is also a solution, and any constant multiple of any solution is also a solution.
Definition
A partial differential equation is an equation that involves an unknown function of variables and (some of) its partial derivatives. That is, for the unknown function
of variables belonging to the open subset of , the -order partial differential equation is defined as
where
and is the partial derivative operator.
Notation
When writing PDEs, it is common to denote partial derivatives using subscripts. For example:
In the general situation that is a function of variables, then denotes the first partial derivative relative to the -th input, denotes the second partial derivative relative to the -th and -th inputs, and so on.
The Greek letter denotes the Laplace operator; if is a function of variables, then
In the physics literature, the Laplace operator is often denoted by ; in the mathematics literature, may also denote the Hessian matrix of .
Classification
Linear and nonlinear equations
A PDE is called linear if it is linear in the unknown and its derivatives. For example, for a function of and , a second order linear PDE is of the form
where and are functions of the independent variables and only. (Often the mixed-partial derivatives and will be equated, but this is not required for the discussion of linearity.)
If the are constants (independent of and ) then the PDE is called linear with constant coefficients. If is zero everywhere then the linear PDE is homogeneous, otherwise it is inhomogeneous. (This is separate from asymptotic homogenization, which studies the effects of high-frequency oscillations in the coefficients upon solutions to PDEs.)
Nearest to linear PDEs are semi-linear PDEs, where only the highest order derivatives appear as linear terms, with coefficients that are functions of the independent variables. The lower order derivatives and the unknown function may appear arbitrarily. For example, a general second order semi-linear PDE in two variables is
In a quasilinear PDE the highest order derivatives likewise appear only as linear terms, but with coefficients possibly functions of the unknown and lower-order derivatives:
Many of the fundamental PDEs in physics are quasilinear, such as the Einstein equations of general relativity and the Navier–Stokes equations describing fluid motion.
A PDE without any linearity properties is called fully nonlinear, and possesses nonlinearities on one or more of the highest-order derivatives. An example is the Monge–Ampère equation, which arises in differential geometry.
Second order equations
The elliptic/parabolic/hyperbolic classification provides a guide to appropriate initial- and boundary conditions and to the smoothness of the solutions. Assuming , the general linear second-order PDE in two independent variables has the form
where the coefficients , , ... may depend upon and . If over a region of the -plane, the PDE is second-order in that region. This form is analogous to the equation for a conic section:
More precisely, replacing by , and likewise for other variables (formally this is done by a Fourier transform), converts a constant-coefficient PDE into a polynomial of the same degree, with the terms of the highest degree (a homogeneous polynomial, here a quadratic form) being most significant for the classification.
Just as one classifies conic sections and quadratic forms into parabolic, hyperbolic, and elliptic based on the discriminant , the same can be done for a second-order PDE at a given point. However, the discriminant in a PDE is given by due to the convention of the term being rather than ; formally, the discriminant (of the associated quadratic form) is , with the factor of 4 dropped for simplicity.
(elliptic partial differential equation): Solutions of elliptic PDEs are as smooth as the coefficients allow, within the interior of the region where the equation and solutions are defined. For example, solutions of Laplace's equation are analytic within the domain where they are defined, but solutions may assume boundary values that are not smooth. The motion of a fluid at subsonic speeds can be approximated with elliptic PDEs, and the Euler–Tricomi equation is elliptic where . By change of variables, the equation can always be expressed in the form: where x and y correspond to changed variables. This justifies Laplace equation as an example of this type.
(parabolic partial differential equation): Equations that are parabolic at every point can be transformed into a form analogous to the heat equation by a change of independent variables. Solutions smooth out as the transformed time variable increases. The Euler–Tricomi equation has parabolic type on the line where . By change of variables, the equation can always be expressed in the form: where x correspond to changed variables. This justifies heat equation, which are of form , as an example of this type.
(hyperbolic partial differential equation): hyperbolic equations retain any discontinuities of functions or derivatives in the initial data. An example is the wave equation. The motion of a fluid at supersonic speeds can be approximated with hyperbolic PDEs, and the Euler–Tricomi equation is hyperbolic where . By change of variables, the equation can always be expressed in the form: where x and y correspond to changed variables. This justifies wave equation as an example of this type.
If there are independent variables , a general linear partial differential equation of second order has the form
The classification depends upon the signature of the eigenvalues of the coefficient matrix .
Elliptic: the eigenvalues are all positive or all negative.
Parabolic: the eigenvalues are all positive or all negative, except one that is zero.
Hyperbolic: there is only one negative eigenvalue and all the rest are positive, or there is only one positive eigenvalue and all the rest are negative.
Ultrahyperbolic: there is more than one positive eigenvalue and more than one negative eigenvalue, and there are no zero eigenvalues.
The theory of elliptic, parabolic, and hyperbolic equations have been studied for centuries, largely centered around or based upon the standard examples of the Laplace equation, the heat equation, and the wave equation.
However, the classification only depends on linearity of the second-order terms and is therefore applicable to semi- and quasilinear PDEs as well. The basic types also extend to hybrids such as the Euler–Tricomi equation; varying from elliptic to hyperbolic for different regions of the domain, as well as higher-order PDEs, but such knowledge is more specialized.
Systems of first-order equations and characteristic surfaces
The classification of partial differential equations can be extended to systems of first-order equations, where the unknown is now a vector with components, and the coefficient matrices are by matrices for . The partial differential equation takes the form
where the coefficient matrices and the vector may depend upon and . If a hypersurface is given in the implicit form
where has a non-zero gradient, then is a characteristic surface for the operator at a given point if the characteristic form vanishes:
The geometric interpretation of this condition is as follows: if data for are prescribed on the surface , then it may be possible to determine the normal derivative of on from the differential equation. If the data on and the differential equation determine the normal derivative of on , then is non-characteristic. If the data on and the differential equation do not determine the normal derivative of on , then the surface is characteristic, and the differential equation restricts the data on : the differential equation is internal to .
A first-order system is elliptic if no surface is characteristic for : the values of on and the differential equation always determine the normal derivative of on .
A first-order system is hyperbolic at a point if there is a spacelike surface with normal at that point. This means that, given any non-trivial vector orthogonal to , and a scalar multiplier , the equation has real roots . The system is strictly hyperbolic if these roots are always distinct. The geometrical interpretation of this condition is as follows: the characteristic form defines a cone (the normal cone) with homogeneous coordinates ζ. In the hyperbolic case, this cone has sheets, and the axis runs inside these sheets: it does not intersect any of them. But when displaced from the origin by η, this axis intersects every sheet. In the elliptic case, the normal cone has no real sheets.
Analytical solutions
Separation of variables
Linear PDEs can be reduced to systems of ordinary differential equations by the important technique of separation of variables. This technique rests on a feature of solutions to differential equations: if one can find any solution that solves the equation and satisfies the boundary conditions, then it is the solution (this also applies to ODEs). We assume as an ansatz that the dependence of a solution on the parameters space and time can be written as a product of terms that each depend on a single parameter, and then see if this can be made to solve the problem.
In the method of separation of variables, one reduces a PDE to a PDE in fewer variables, which is an ordinary differential equation if in one variable – these are in turn easier to solve.
This is possible for simple PDEs, which are called separable partial differential equations, and the domain is generally a rectangle (a product of intervals). Separable PDEs correspond to diagonal matrices – thinking of "the value for fixed " as a coordinate, each coordinate can be understood separately.
This generalizes to the method of characteristics, and is also used in integral transforms.
Method of characteristics
The characteristic surface in dimensional space is called a characteristic curve.
In special cases, one can find characteristic curves on which the first-order PDE reduces to an ODE – changing coordinates in the domain to straighten these curves allows separation of variables, and is called the method of characteristics.
More generally, applying the method to first-order PDEs in higher dimensions, one may find characteristic surfaces.
Integral transform
An integral transform may transform the PDE to a simpler one, in particular, a separable PDE. This corresponds to diagonalizing an operator.
An important example of this is Fourier analysis, which diagonalizes the heat equation using the eigenbasis of sinusoidal waves.
If the domain is finite or periodic, an infinite sum of solutions such as a Fourier series is appropriate, but an integral of solutions such as a Fourier integral is generally required for infinite domains. The solution for a point source for the heat equation given above is an example of the use of a Fourier integral.
Change of variables
Often a PDE can be reduced to a simpler form with a known solution by a suitable change of variables. For example, the Black–Scholes equation
is reducible to the heat equation
by the change of variables
Fundamental solution
Inhomogeneous equations can often be solved (for constant coefficient PDEs, always be solved) by finding the fundamental solution (the solution for a point source ), then taking the convolution with the boundary conditions to get the solution.
This is analogous in signal processing to understanding a filter by its impulse response.
Superposition principle
The superposition principle applies to any linear system, including linear systems of PDEs. A common visualization of this concept is the interaction of two waves in phase being combined to result in a greater amplitude, for example . The same principle can be observed in PDEs where the solutions may be real or complex and additive. If and are solutions of linear PDE in some function space , then with any constants and are also a solution of that PDE in the same function space.
Methods for non-linear equations
There are no generally applicable analytical methods to solve nonlinear PDEs. Still, existence and uniqueness results (such as the Cauchy–Kowalevski theorem) are often possible, as are proofs of important qualitative and quantitative properties of solutions (getting these results is a major part of analysis).
Nevertheless, some techniques can be used for several types of equations. The -principle is the most powerful method to solve underdetermined equations. The Riquier–Janet theory is an effective method for obtaining information about many analytic overdetermined systems.
The method of characteristics can be used in some very special cases to solve nonlinear partial differential equations.
In some cases, a PDE can be solved via perturbation analysis in which the solution is considered to be a correction to an equation with a known solution. Alternatives are numerical analysis techniques from simple finite difference schemes to the more mature multigrid and finite element methods. Many interesting problems in science and engineering are solved in this way using computers, sometimes high performance supercomputers.
Lie group method
From 1870 Sophus Lie's work put the theory of differential equations on a more satisfactory foundation. He showed that the integration theories of the older mathematicians can, by the introduction of what are now called Lie groups, be referred, to a common source; and that ordinary differential equations which admit the same infinitesimal transformations present comparable difficulties of integration. He also emphasized the subject of transformations of contact.
A general approach to solving PDEs uses the symmetry property of differential equations, the continuous infinitesimal transformations of solutions to solutions (Lie theory). Continuous group theory, Lie algebras and differential geometry are used to understand the structure of linear and nonlinear partial differential equations for generating integrable equations, to find its Lax pairs, recursion operators, Bäcklund transform and finally finding exact analytic solutions to the PDE.
Symmetry methods have been recognized to study differential equations arising in mathematics, physics, engineering, and many other disciplines.
Semi-analytical methods
The Adomian decomposition method, the Lyapunov artificial small parameter method, and his homotopy perturbation method are all special cases of the more general homotopy analysis method. These are series expansion methods, and except for the Lyapunov method, are independent of small physical parameters as compared to the well known perturbation theory, thus giving these methods greater flexibility and solution generality.
Numerical solutions
The three most widely used numerical methods to solve PDEs are the finite element method (FEM), finite volume methods (FVM) and finite difference methods (FDM), as well other kind of methods called meshfree methods, which were made to solve problems where the aforementioned methods are limited. The FEM has a prominent position among these methods and especially its exceptionally efficient higher-order version hp-FEM. Other hybrid versions of FEM and Meshfree methods include the generalized finite element method (GFEM), extended finite element method (XFEM), spectral finite element method (SFEM), meshfree finite element method, discontinuous Galerkin finite element method (DGFEM), element-free Galerkin method (EFGM), interpolating element-free Galerkin method (IEFGM), etc.
Finite element method
The finite element method (FEM) (its practical application often known as finite element analysis (FEA)) is a numerical technique for finding approximate solutions of partial differential equations (PDE) as well as of integral equations. The solution approach is based either on eliminating the differential equation completely (steady state problems), or rendering the PDE into an approximating system of ordinary differential equations, which are then numerically integrated using standard techniques such as Euler's method, Runge–Kutta, etc.
Finite difference method
Finite-difference methods are numerical methods for approximating the solutions to differential equations using finite difference equations to approximate derivatives.
Finite volume method
Similar to the finite difference method or finite element method, values are calculated at discrete places on a meshed geometry. "Finite volume" refers to the small volume surrounding each node point on a mesh. In the finite volume method, surface integrals in a partial differential equation that contain a divergence term are converted to volume integrals, using the divergence theorem. These terms are then evaluated as fluxes at the surfaces of each finite volume. Because the flux entering a given volume is identical to that leaving the adjacent volume, these methods conserve mass by design.
Neural networks
Weak solutions
Weak solutions are functions that satisfy the PDE, yet in other meanings than regular sense. The meaning for this term may differ with context, and one of the most commonly used definitions is based on the notion of distributions.
An example for the definition of a weak solution is as follows:
Consider the boundary-value problem given by:
where denotes a second-order partial differential operator in divergence form.
We say a is a weak solution if
for every , which can be derived by a formal integral by parts.
An example for a weak solution is as follows:
is a weak solution satisfying
in distributional sense, as formally,
Theoretical Studies
As a branch of pure mathematics, the theoretical studies of PDEs focus on the criteria for a solution to exist, the properties of a solution, and finding its formula is often secondary.
Well-posedness
Well-posedness refers to a common schematic package of information about a PDE. To say that a PDE is well-posed, one must have:
an existence and uniqueness theorem, asserting that by the prescription of some freely chosen functions, one can single out one specific solution of the PDE
by continuously changing the free choices, one continuously changes the corresponding solution
This is, by the necessity of being applicable to several different PDE, somewhat vague. The requirement of "continuity", in particular, is ambiguous, since there are usually many inequivalent means by which it can be rigorously defined. It is, however, somewhat unusual to study a PDE without specifying a way in which it is well-posed.
Regularity
Regularity refers to the integrability and differentiability of weak solutions, which can often be represented by Sobolev spaces.
This problem arise due to the difficulty in searching for classical solutions. Researchers often tend to find weak solutions at first and then find out whether it is smooth enough to be qualified as a classical solution.
Results from functional analysis are often used in this field of study.
| Mathematics | Calculus and analysis | null |
52565 | https://en.wikipedia.org/wiki/Partial%20derivative | Partial derivative | In mathematics, a partial derivative of a function of several variables is its derivative with respect to one of those variables, with the others held constant (as opposed to the total derivative, in which all variables are allowed to vary). Partial derivatives are used in vector calculus and differential geometry.
The partial derivative of a function with respect to the variable is variously denoted by
It can be thought of as the rate of change of the function in the -direction.
Sometimes, for the partial derivative of with respect to is denoted as Since a partial derivative generally has the same arguments as the original function, its functional dependence is sometimes explicitly signified by the notation, such as in:
The symbol used to denote partial derivatives is ∂. One of the first known uses of this symbol in mathematics is by Marquis de Condorcet from 1770, who used it for partial differences. The modern partial derivative notation was created by Adrien-Marie Legendre (1786), although he later abandoned it; Carl Gustav Jacob Jacobi reintroduced the symbol in 1841.
Definition
Like ordinary derivatives, the partial derivative is defined as a limit. Let be an open subset of and a function. The partial derivative of at the point with respect to the -th variable is defined as
Where is the unit vector of -th variable . Even if all partial derivatives exist at a given point , the function need not be continuous there. However, if all partial derivatives exist in a neighborhood of and are continuous there, then is totally differentiable in that neighborhood and the total derivative is continuous. In this case, it is said that is a function. This can be used to generalize for vector valued functions, by carefully using a componentwise argument.
The partial derivative can be seen as another function defined on and can again be partially differentiated. If the direction of derivative is repeated, it is called a mixed partial derivative. If all mixed second order partial derivatives are continuous at a point (or on a set), is termed a function at that point (or on that set); in this case, the partial derivatives can be exchanged by Clairaut's theorem:
Notation
For the following examples, let be a function in , , and .
First-order partial derivatives:
Second-order partial derivatives:
Second-order mixed derivatives:
Higher-order partial and mixed derivatives:
When dealing with functions of multiple variables, some of these variables may be related to each other, thus it may be necessary to specify explicitly which variables are being held constant to avoid ambiguity. In fields such as statistical mechanics, the partial derivative of with respect to , holding and constant, is often expressed as
Conventionally, for clarity and simplicity of notation, the partial derivative function and the value of the function at a specific point are conflated by including the function arguments when the partial derivative symbol (Leibniz notation) is used. Thus, an expression like
is used for the function, while
might be used for the value of the function at the point However, this convention breaks down when we want to evaluate the partial derivative at a point like In such a case, evaluation of the function must be expressed in an unwieldy manner as
or
in order to use the Leibniz notation. Thus, in these cases, it may be preferable to use the Euler differential operator notation with as the partial derivative symbol with respect to the -th variable. For instance, one would write for the example described above, while the expression represents the partial derivative function with respect to the first variable.
For higher order partial derivatives, the partial derivative (function) of with respect to the -th variable is denoted That is, so that the variables are listed in the order in which the derivatives are taken, and thus, in reverse order of how the composition of operators is usually notated. Of course, Clairaut's theorem implies that as long as comparatively mild regularity conditions on are satisfied.
Gradient
An important example of a function of several variables is the case of a scalar-valued function on a domain in Euclidean space (e.g., on or In this case has a partial derivative with respect to each variable . At the point , these partial derivatives define the vector
This vector is called the gradient of at . If is differentiable at every point in some domain, then the gradient is a vector-valued function which takes the point to the vector . Consequently, the gradient produces a vector field.
A common abuse of notation is to define the del operator () as follows in three-dimensional Euclidean space with unit vectors
Or, more generally, for -dimensional Euclidean space with coordinates and unit vectors
Directional derivative
Example
Suppose that is a function of more than one variable. For instance,
The graph of this function defines a surface in Euclidean space. To every point on this surface, there are an infinite number of tangent lines. Partial differentiation is the act of choosing one of these lines and finding its slope. Usually, the lines of most interest are those that are parallel to the -plane, and those that are parallel to the -plane (which result from holding either or constant, respectively).
To find the slope of the line tangent to the function at and parallel to the -plane, we treat as a constant. The graph and this plane are shown on the right. Below, we see how the function looks on the plane . By finding the derivative of the equation while assuming that is a constant, we find that the slope of at the point is:
So at , by substitution, the slope is . Therefore,
at the point . That is, the partial derivative of with respect to at is , as shown in the graph.
The function can be reinterpreted as a family of functions of one variable indexed by the other variables:
In other words, every value of defines a function, denoted , which is a function of one variable . That is,
In this section the subscript notation denotes a function contingent on a fixed value of , and not a partial derivative.
Once a value of is chosen, say , then determines a function which traces a curve on the -plane:
In this expression, is a , not a , so is a function of only one real variable, that being . Consequently, the definition of the derivative for a function of one variable applies:
The above procedure can be performed for any choice of . Assembling the derivatives together into a function gives a function which describes the variation of in the direction:
This is the partial derivative of with respect to . Here '' is a rounded 'd' called the partial derivative symbol; to distinguish it from the letter 'd', '' is sometimes pronounced "partial".
Higher order partial derivatives
Second and higher order partial derivatives are defined analogously to the higher order derivatives of univariate functions. For the function the "own" second partial derivative with respect to is simply the partial derivative of the partial derivative (both with respect to ):
The cross partial derivative with respect to and is obtained by taking the partial derivative of with respect to , and then taking the partial derivative of the result with respect to , to obtain
Schwarz's theorem states that if the second derivatives are continuous, the expression for the cross partial derivative is unaffected by which variable the partial derivative is taken with respect to first and which is taken second. That is,
or equivalently
Own and cross partial derivatives appear in the Hessian matrix which is used in the second order conditions in optimization problems.
The higher order partial derivatives can be obtained by successive differentiation
Antiderivative analogue
There is a concept for partial derivatives that is analogous to antiderivatives for regular derivatives. Given a partial derivative, it allows for the partial recovery of the original function.
Consider the example of
The so-called partial integral can be taken with respect to (treating as constant, in a similar manner to partial differentiation):
Here, the constant of integration is no longer a constant, but instead a function of all the variables of the original function except . The reason for this is that all the other variables are treated as constant when taking the partial derivative, so any function which does not involve will disappear when taking the partial derivative, and we have to account for this when we take the antiderivative. The most general way to represent this is to have the constant represent an unknown function of all the other variables.
Thus the set of functions where is any one-argument function, represents the entire set of functions in variables that could have produced the -partial derivative
If all the partial derivatives of a function are known (for example, with the gradient), then the antiderivatives can be matched via the above process to reconstruct the original function up to a constant. Unlike in the single-variable case, however, not every set of functions can be the set of all (first) partial derivatives of a single function. In other words, not every vector field is conservative.
Applications
Geometry
The volume of a cone depends on the cone's height and its radius according to the formula
The partial derivative of with respect to is
which represents the rate with which a cone's volume changes if its radius is varied and its height is kept constant. The partial derivative with respect to equals which represents the rate with which the volume changes if its height is varied and its radius is kept constant.
By contrast, the total derivative of with respect to and are respectively
The difference between the total and partial derivative is the elimination of indirect dependencies between variables in partial derivatives.
If (for some arbitrary reason) the cone's proportions have to stay the same, and the height and radius are in a fixed ratio ,
This gives the total derivative with respect to ,
which simplifies to
Similarly, the total derivative with respect to is
The total derivative with respect to and of the volume intended as scalar function of these two variables is given by the gradient vector
Optimization
Partial derivatives appear in any calculus-based optimization problem with more than one choice variable. For example, in economics a firm may wish to maximize profit with respect to the choice of the quantities and of two different types of output. The first order conditions for this optimization are . Since both partial derivatives and will generally themselves be functions of both arguments and , these two first order conditions form a system of two equations in two unknowns.
Thermodynamics, quantum mechanics and mathematical physics
Partial derivatives appear in thermodynamic equations like Gibbs-Duhem equation, in quantum mechanics as in Schrödinger wave equation, as well as in other equations from mathematical physics. The variables being held constant in partial derivatives here can be ratios of simple variables like mole fractions in the following example involving the Gibbs energies in a ternary mixture system:
Express mole fractions of a component as functions of other components' mole fraction and binary mole ratios:
Differential quotients can be formed at constant ratios like those above:
Ratios X, Y, Z of mole fractions can be written for ternary and multicomponent systems:
which can be used for solving partial differential equations like:
This equality can be rearranged to have differential quotient of mole fractions on one side.
Image resizing
Partial derivatives are key to target-aware image resizing algorithms. Widely known as seam carving, these algorithms require each pixel in an image to be assigned a numerical 'energy' to describe their dissimilarity against orthogonal adjacent pixels. The algorithm then progressively removes rows or columns with the lowest energy. The formula established to determine a pixel's energy (magnitude of gradient at a pixel) depends heavily on the constructs of partial derivatives.
Economics
Partial derivatives play a prominent role in economics, in which most functions describing economic behaviour posit that the behaviour depends on more than one variable. For example, a societal consumption function may describe the amount spent on consumer goods as depending on both income and wealth; the marginal propensity to consume is then the partial derivative of the consumption function with respect to income.
| Mathematics | Calculus and analysis | null |
52636 | https://en.wikipedia.org/wiki/Boiling | Boiling | Boiling or ebullition is the rapid phase transition from liquid to gas or vapour; the reverse of boiling is condensation. Boiling occurs when a liquid is heated to its boiling point, so that the vapour pressure of the liquid is equal to the pressure exerted on the liquid by the surrounding atmosphere. Boiling and evaporation are the two main forms of liquid vapourization.
There are two main types of boiling: nucleate boiling where small bubbles of vapour form at discrete points, and critical heat flux boiling where the boiling surface is heated above a certain critical temperature and a film of vapour forms on the surface. Transition boiling is an intermediate, unstable form of boiling with elements of both types. The boiling point of water is 100 °C or 212 °F but is lower with the decreased atmospheric pressure found at higher altitudes.
Boiling water is used as a method of making it potable by killing microbes and viruses that may be present. The sensitivity of different micro-organisms to heat varies, but if water is held at for one minute, most micro-organisms and viruses are inactivated. Ten minutes at a temperature of 70 °C (158 °F) is also sufficient to inactivate most bacteria.
Boiling water is also used in several cooking methods including boiling, steaming, and poaching.
Types
Free convection
The lowest heat flux seen in boiling is only sufficient to cause [natural convection], where the warmer fluid rises due to its slightly lower density. This condition occurs only when the superheat is very low, meaning that the hot surface near the fluid is nearly the same temperature as the boiling point.
Nucleate
Nucleate boiling is characterised by the growth of bubbles or pops on a heated surface (heterogeneous nucleation), which rises from discrete points on a surface, whose temperature is only slightly above the temperature of the liquid. In general, the number of nucleation sites is increased by an increasing surface temperature.
An irregular surface of the boiling vessel (i.e., increased surface roughness) or additives to the fluid (i.e., surfactants and/or nanoparticles) facilitate nucleate boiling over a broader temperature range, while an exceptionally smooth surface, such as plastic, lends itself to superheating. Under these conditions, a heated liquid may show boiling delay and the temperature may go somewhat above the boiling point without boiling.
Homogeneous nucleation, where the bubbles form from the surrounding liquid instead of on a surface, can occur if the liquid is warmer in its center, and cooler at the surfaces of the container. This can be done, for instance, in a microwave oven, which heats the water and not the container.
Critical heat flux
Critical heat flux (CHF) describes the thermal limit of a phenomenon where a phase change occurs during heating (such as bubbles forming on a metal surface used to heat water), which suddenly decreases the efficiency of heat transfer, thus causing localised overheating of the heating surface. As the boiling surface is heated above a critical temperature, a film of vapour forms on the surface. Since this vapour film is much less capable of carrying heat away from the surface, the temperature rises very rapidly beyond this point into the transition boiling regime. The point at which this occurs is dependent on the characteristics of boiling fluid and the heating surface in question.
Transition
Transition boiling may be defined as the unstable boiling, which occurs at surface temperatures between the maximum attainable in nucleate and the minimum attainable in film boiling.
The formation of bubbles in a heated liquid is a complex physical process which often involves cavitation and acoustic effects, such as the broad-spectrum hiss one hears in a kettle not yet heated to the point where bubbles boil to the surface.
Film
If a surface heating the liquid is significantly hotter than the liquid then film boiling will occur, where a thin layer of vapour, which has low thermal conductivity, insulates the surface. This condition of a vapour film insulating the surface from the liquid characterises film boiling.
Influence of geometry
Pool boiling
"Pool boiling" refers to boiling where there is no forced convective flow. Instead, the flow occurs due to density gradients. It can experience any of the regimes mentioned above.
Flow boiling
"Flow boiling" occurs when the boiling fluid circulates, typically through pipes. Its movement can be powered by pumps, such as in power plants, or by density gradients, such as in a thermosiphon or a heat pipe. Flows in flow boiling are often characterised by a void fraction parameter, which indicates the fraction of the volume in the system that is vapor. One can use this fraction and the densities to calculate the vapor quality, which refers to the mass fraction that is in the gas phase. Flow boiling can be very complex, with heavy influences of density, flow rates, and heat flux, as well as surface tension. The same system may have regions that are liquid, gas, and two-phase flow. Such two phase regimes can lead to some of the best heat transfer coefficients of any system.
Confined boiling
Confined boiling refers to boiling in confined geometries, typically characterized by a Bond number that compares the gap spacing to the capillary length. Confined boiling regimes begin to play a major role when Bo < 0.5. This boiling regime is dominated by "vapour stem bubbles" left behind after vapour departs. These bubbles act as seeds for vapor growth. Confined boiling typically has higher heat transfer coefficient but a lower CHF than pool boiling. CHF occurs when the vapor momentum force at the two-phase interface balances the combined surface tension and hydrostatic forces, leading to irreversible growth of the dry spot. Confined boiling is particularly promising for electronics cooling.
Physics
The boiling point of an element at a given pressure is a characteristic attribute of the element. This is also true for many simple compounds including water and simple alcohols. Once boiling has started and provided that boiling remains stable and the pressure is constant, the temperature of the boiling liquid remains constant. This attribute led to the adoption of boiling points as the definition of 100 °C.
Distillation
Mixtures of volatile liquids have a boiling point specific to that mixture producing vapour with a constant mix of components - the constant boiling mixture. This attribute allows mixtures of liquids to be separated or partly separated by boiling and is best known as a means of separating ethanol from water.
Uses
Refrigeration and air conditioning
Most types of refrigeration and some type of air-conditioning work by compressing a gas so that it becomes liquid and then allowing it to boil. This adsorbs heat from the surroundings cooling the fridge or freezer or cooling the air entering a building. Typical liquids include propane, ammonia, carbon dioxide or nitrogen.
For making water potable
As a method of disinfecting water, bringing it to its boiling point at , is the oldest and most effective way since it does not affect the taste, it is effective despite contaminants or particles present in it, and is a single step process which eliminates most microbes responsible for causing intestine related diseases. The boiling point of water is at sea level and at normal barometric pressure. In places having a proper water purification system, it is recommended only as an emergency treatment method or for obtaining potable water in the wilderness or in rural areas, as it cannot remove chemical toxins or impurities.
The elimination of micro-organisms by boiling follows first-order kinetics—at high temperatures, it is achieved in less time and at lower temperatures, in more time. The heat sensitivity of micro-organisms varies, at , Giardia species (which cause giardiasis) can take ten minutes for complete inactivation, most intestine affecting microbes and E. coli (gastroenteritis) take less than a minute; at boiling point, Vibrio cholerae (cholera) takes ten seconds and hepatitis A virus (causes the symptom of jaundice), one minute. Boiling does not ensure the elimination of all micro-organisms; the bacterial spores Clostridium can survive at but are not water-borne or intestine affecting. Thus for human health, complete sterilization of water is not required.
The traditional advice of boiling water for ten minutes is mainly for additional safety, since microbes start getting eliminated at temperatures greater than and bringing it to its boiling point is also a useful indication that can be seen without the help of a thermometer, and by this time, the water is disinfected. Though the boiling point decreases with increasing altitude, it is not enough to affect the disinfecting process.
In cooking
Boiling is the method of cooking food in boiling water or other water-based liquids such as stock or milk. Simmering is gentle boiling, while in poaching the cooking liquid moves but scarcely bubbles.
The boiling point of water is typically considered to be , especially at sea level. Pressure and a change in the composition of the liquid may alter the boiling point of the liquid. High elevation cooking generally takes longer since boiling point is a function of atmospheric pressure. At an elevation of about , water boils at approximately . Depending on the type of food and the elevation, the boiling water may not be hot enough to cook the food properly. Similarly, increasing the pressure as in a pressure cooker raises the temperature of the contents above the open air boiling point.
Boil-in-the-bag
Also known as "boil-in-bag", this involves heating or cooking ready-made foods sealed in a thick plastic bag. The bag containing the food, often frozen, is submerged in boiling water for a prescribed time. The resulting dishes can be prepared with greater convenience as no pots or pans are dirtied in the process. Such meals are available for camping as well as home dining.
Contrast with evaporation
At any given temperature, the molecules in a liquid have varying kinetic energies. Some high energy particles on the liquid surface may have enough energy to escape the intermolecular forces of attraction of the liquid and become a gas. This is called evaporation.
Evaporation only happens on the surface while boiling happens throughout the liquid.
When a liquid reaches its boiling point bubbles of gas form in it which rise into the surface and burst into the air. This process is called boiling. If the boiling liquid is heated more strongly the temperature does not rise but the liquid boils more quickly.
This distinction is exclusive to the liquid-to-gas transition; any transition directly from solid to gas is always referred to as sublimation regardless of whether it is at its boiling point or not.
| Physical sciences | Phase transitions | null |
52644 | https://en.wikipedia.org/wiki/Cysteine | Cysteine | Cysteine (; symbol Cys or C) is a semiessential proteinogenic amino acid with the formula . The thiol side chain in cysteine enables the formation of disulfide bonds, and often participates in enzymatic reactions as a nucleophile. Cysteine is chiral, but both D and L-cysteine are found in nature. LCysteine is a protein monomer in all biota, and D-cysteine acts as a signaling molecule in mammalian nervous systems. Cysteine is named after its discovery in urine, which comes from the urinary bladder or cyst, from Greek κύστις kýstis, "bladder".
The thiol is susceptible to oxidation to give the disulfide derivative cystine, which serves an important structural role in many proteins. In this case, the symbol Cyx is sometimes used. The deprotonated form can generally be described by the symbol Cym as well.
When used as a food additive, cysteine has the E number E920.
Cysteine is encoded by the codons UGU and UGC.
Structure
Like other amino acids (not as a residue of a protein), cysteine exists as a zwitterion. Cysteine has chirality in the older / notation based on homology to - and -glyceraldehyde. In the newer R/S system of designating chirality, based on the atomic numbers of atoms near the asymmetric carbon, cysteine (and selenocysteine) have R chirality, because of the presence of sulfur (or selenium) as a second neighbor to the asymmetric carbon atom. The remaining chiral amino acids, having lighter atoms in that position, have S chirality. Replacing sulfur with selenium gives selenocysteine.
Dietary sources
Cysteinyl is a residue in high-protein foods. Some foods considered rich in cysteine include poultry, eggs, beef, and whole grains. In high-protein diets, cysteine may be partially responsible for reduced blood pressure and stroke risk. Although classified as a nonessential amino acid, in rare cases, cysteine may be essential for infants, the elderly, and individuals with certain metabolic diseases or who suffer from malabsorption syndromes. Cysteine can usually be synthesized by the human body under normal physiological conditions if a sufficient quantity of methionine is available.
Industrial sources
The majority of -cysteine is obtained industrially by hydrolysis of animal materials, such as poultry feathers or hog hair. Despite widespread rumor, human hair is rarely a source material. Indeed, food additive or cosmetic product manufactures may not legally source from human hair in the European Union.
Some animal-originating sources of -cysteine as a food additive contravene kosher, halal, vegan, or vegetarian diets. To avoid this problem, synthetic -cysteine, compliant with Jewish kosher and Muslim halal laws, is also available, albeit at a higher price. The typical synthetic route involves fermentation with an artificial E. coli strain.
Alternatively, Evonik (formerly Degussa) introduced a route from substituted thiazolines. Pseudomonas thiazolinophilum hydrolyzes racemic 2amino-Δ2thiazoline-4carboxylic acid to cysteine.
Biosynthesis
In animals, biosynthesis begins with the amino acid serine. The sulfur is derived from methionine, which is converted to homocysteine through the intermediate S-adenosylmethionine. Cystathionine beta-synthase then combines homocysteine and serine to form the asymmetrical thioether cystathionine. The enzyme cystathionine gamma-lyase converts the cystathionine into cysteine and alpha-ketobutyrate. In plants and bacteria, cysteine biosynthesis also starts from serine, which is converted to O-acetylserine by the enzyme serine transacetylase. The enzyme cysteine synthase, using sulfide sources, converts this ester into cysteine, releasing acetate.
Biological functions
The cysteine sulfhydryl group is nucleophilic and easily oxidized. The reactivity is enhanced when the thiol is ionized, and cysteine residues in proteins have pKa values close to neutrality, so are often in their reactive thiolate form in the cell. Because of its high reactivity, the sulfhydryl group of cysteine has numerous biological functions.
Precursor to the antioxidant glutathione
Due to the ability of thiols to undergo redox reactions, cysteine and cysteinyl residues have antioxidant properties. Its antioxidant properties are typically expressed in the tripeptide glutathione, which occurs in humans and other organisms. The systemic availability of oral glutathione (GSH) is negligible; so it must be biosynthesized from its constituent amino acids, cysteine, glycine, and glutamic acid. While glutamic acid is usually sufficient because amino acid nitrogen is recycled through glutamate as an intermediary, dietary cysteine and glycine supplementation can improve synthesis of glutathione.
Precursor to iron-sulfur clusters
Cysteine is an important source of sulfide in human metabolism. The sulfide in iron-sulfur clusters and in nitrogenase is extracted from cysteine, which is converted to alanine in the process.
Metal ion binding
Beyond the iron-sulfur proteins, many other metal cofactors in enzymes are bound to the thiolate substituent of cysteinyl residues. Examples include zinc in zinc fingers and alcohol dehydrogenase, copper in the blue copper proteins, iron in cytochrome P450, and nickel in the [NiFe]-hydrogenases. The sulfhydryl group also has a high affinity for heavy metals, so that proteins containing cysteine, such as metallothionein, will bind metals such as mercury, lead, and cadmium tightly.
Roles in protein structure
In the translation of messenger RNA molecules to produce polypeptides, cysteine is coded for by the UGU and UGC codons.
Cysteine has traditionally been considered to be a hydrophilic amino acid, based largely on the chemical parallel between its sulfhydryl group and the hydroxyl groups in the side chains of other polar amino acids. However, the cysteine side chain has been shown to stabilize hydrophobic interactions in micelles to a greater degree than the side chain in the nonpolar amino acid glycine and the polar amino acid serine. In a statistical analysis of the frequency with which amino acids appear in various proteins, cysteine residues were found to associate with hydrophobic regions of proteins. Their hydrophobic tendency was equivalent to that of known nonpolar amino acids such as methionine and tyrosine (tyrosine is polar aromatic but also hydrophobic), those of which were much greater than that of known polar amino acids such as serine and threonine. Hydrophobicity scales, which rank amino acids from most hydrophobic to most hydrophilic, consistently place cysteine towards the hydrophobic end of the spectrum, even when they are based on methods that are not influenced by the tendency of cysteines to form disulfide bonds in proteins. Therefore, cysteine is now often grouped among the hydrophobic amino acids, though it is sometimes also classified as slightly polar, or polar.
Most cysteine residues are covalently bonded to other cysteine residues to form disulfide bonds, which play an important role in the folding and stability of some proteins, usually proteins secreted to the extracellular medium. Since most cellular compartments are reducing environments, disulfide bonds are generally unstable in the cytosol with some exceptions as noted below.
Disulfide bonds in proteins are formed by oxidation of the sulfhydryl group of cysteine residues. The other sulfur-containing amino acid, methionine, cannot form disulfide bonds. More aggressive oxidants convert cysteine to the corresponding sulfinic acid and sulfonic acid. Cysteine residues play a valuable role by crosslinking proteins, which increases the rigidity of proteins and also functions to confer proteolytic resistance (since protein export is a costly process, minimizing its necessity is advantageous). Inside the cell, disulfide bridges between cysteine residues within a polypeptide support the protein's tertiary structure. Insulin is an example of a protein with cystine crosslinking, wherein two separate peptide chains are connected by a pair of disulfide bonds.
Protein disulfide isomerases catalyze the proper formation of disulfide bonds; the cell transfers dehydroascorbic acid to the endoplasmic reticulum, which oxidizes the environment. In this environment, cysteines are, in general, oxidized to cystine and are no longer functional as a nucleophiles.
Aside from its oxidation to cystine, cysteine participates in numerous post-translational modifications. The nucleophilic sulfhydryl group allows cysteine to conjugate to other groups, e.g., in prenylation. Ubiquitin ligases transfer ubiquitin to its pendant, proteins, and caspases, which engage in proteolysis in the apoptotic cycle. Inteins often function with the help of a catalytic cysteine. These roles are typically limited to the intracellular milieu, where the environment is reducing, and cysteine is not oxidized to cystine.
Evolutionary role of cysteine
Cysteine is considered a "newcomer" amino acid, being the 17th amino acid incorporated into the genetic code. Similar to other later-added amino acids such as methionine, tyrosine, and tryptophan, cysteine exhibits strong nucleophilic and redox-active properties. These properties contribute to the depletion of cysteine from respiratory chain complexes, such as Complexes I and IV, since reactive oxygen species (ROS) produced by the respiratory chain can react with the cysteine residues in these complexes, leading to dysfunctional proteins and potentially contributing to aging. The primary response of a protein to ROS is the oxidation of cysteine and the loss of free thiol groups, resulting in increased thiyl radicals and associated protein cross-linking. In contrast, another sulfur-containing, redox-active amino acid, methionine, does not exhibit these biochemical properties and its content is relatively upregulated in mitochondrially encoded proteins.
Applications
Cysteine, mainly the -enantiomer, is a precursor in the food, pharmaceutical, and personal-care industries. One of the largest applications is the production of flavors. For example, the reaction of cysteine with sugars in a Maillard reaction yields meat flavors. -Cysteine is also used as a processing aid for baking.
In the field of personal care, cysteine is used for permanent-wave applications, predominantly in Asia. Again, the cysteine is used for breaking up the disulfide bonds in the hair's keratin.
Cysteine is a very popular target for site-directed labeling experiments to investigate biomolecular structure and dynamics. Maleimides selectively attach to cysteine using a covalent Michael addition. Site-directed spin labeling for EPR or paramagnetic relaxation-enhanced NMR also uses cysteine extensively.
Reducing toxic effects of alcohol
Cysteine has been proposed as a preventive or antidote for some of the negative effects of alcohol, including liver damage and hangover. It counteracts the poisonous effects of acetaldehyde. It binds to acetaldehyde to form the low-toxicity heterocycle methylthioproline.
In a rat study, test animals received an LD90 dose of acetaldehyde. Those that received cysteine had an 80% survival rate; when both cysteine and thiamine were administered, all animals survived. The control group had a 10% survival rate.
In 2020 an article was published that suggests L-cysteine might also work in humans.
N-Acetylcysteine
N-Acetyl--cysteine is a derivative of cysteine wherein an acetyl group is attached to the nitrogen atom. This compound is sold as a dietary supplement, and used as an antidote in cases of acetaminophen overdose.
Sheep
Cysteine is required by sheep to produce wool. It is an essential amino acid that is taken in from their feed. As a consequence, during drought conditions, sheep produce less wool; however, transgenic sheep that can make their own cysteine have been developed.
Chemical reactions
Being multifunctional, cysteine undergoes a variety of reactions. Much attention has focused on protecting the sulfhydryl group. Methylation of cysteine gives S-methylcysteine. Treatment with formaldehyde gives the thiazolidine thioproline. Cysteine forms a variety of coordination complexes upon treatment with metal ions.
Safety
Relative to most other amino acids, cysteine is much more toxic.
History
In 1884 German chemist Eugen Baumann found that when cystine was treated with a reducing agent, cystine revealed itself to be a dimer of a monomer which he named "cysteïne".
| Biology and health sciences | Amino acids | Biology |
52648 | https://en.wikipedia.org/wiki/Camera | Camera | A camera is an instrument used to capture and store images and videos, either digitally via an electronic image sensor, or chemically via a light-sensitive material such as photographic film. As a pivotal technology in the fields of photography and videography, cameras have played a significant role in the progression of visual arts, media, entertainment, surveillance, and scientific research. The invention of the camera dates back to the 19th century and has since evolved with advancements in technology, leading to a vast array of types and models in the 21st century.
Cameras function through a combination of multiple mechanical components and principles. These include exposure control, which regulates the amount of light reaching the sensor or film; the lens, which focuses the light; the viewfinder, which allows the user to preview the scene; and the film or sensor, which captures the image.
Several types of cameras exist, each suited to specific uses and offering unique capabilities. Single-lens reflex (SLR) cameras provide real-time, exact imaging through the lens. Large-format and medium-format cameras offer higher image resolution and are often used in professional and artistic photography. Compact cameras, known for their portability and simplicity, are popular in consumer photography. Rangefinder cameras, with separate viewing and imaging systems, were historically widely used in photojournalism. Motion picture cameras are specialized for filming cinematic content, while digital cameras, which became prevalent in the late 20th and early 21st century, use electronic sensors to capture and store images.
The rapid development of smartphone camera technology in the 21st century has blurred the lines between dedicated cameras and multifunctional devices, profoundly influencing how society creates, shares, and consumes visual content.
History
19th century
Beginning with the use of the camera obscura and transitioning to complex photographic cameras, the evolution of the technology in the 19th century was driven by pioneers like Thomas Wedgwood, Nicéphore Niépce, and Henry Fox Talbot. First using the camera obscura for chemical experiments, they ultimately created cameras specifically for chemical photography, and later reduced the camera's size and optimized lens configurations.
The introduction of the daguerreotype process in 1839 facilitated commercial camera manufacturing, with various producers contributing diverse designs. As camera manufacturing became a specialized trade in the 1850s, designs and sizes were standardized.
The latter half of the century witnessed the advent of dry plates and roll-film, prompting a shift towards smaller and more cost-effective cameras, epitomized by the original Kodak camera, first produced in 1888. This period also saw significant advancements in lens technology and the emergence of color photography, leading to a surge in camera ownership.
20th century
The first half of the 20th century saw continued miniaturization and the integration of new manufacturing materials. After World War I, Germany took the lead in camera development, spearheading industry consolidation and producing precision-made cameras. The industry saw significant product launches such as the Leica camera and the Contax, which were enabled by advancements in film and lens designs. Additionally, there was a marked increase in accessibility to cinematography for amateurs with Eastman Kodak's production of the first 16-mm and 8-mm reversal safety films. The World War II era saw a focus on the development of specialized aerial reconnaissance and instrument-recording equipment, even as the overall pace of non-military camera innovation slowed.
In the second half of the century, Japanese manufacturers in particular advanced camera technology. From the introduction of the affordable Ricohflex III TLR in 1952 to the first 35mm SLR with automatic exposure, the Olympus AutoEye in 1960, new designs and features continuously emerged. Electronics became integral to camera design in the 1970s, evident in models like Polaroid's SX-70 and Canon's AE-1.
Transition to digital photography marked the late 20th century, culminating in digital camera sales surpassing film cameras in the United States by 2003. In contrast, the film camera industry in the UK, Western Europe, and the USA declined during this period, while manufacturing continued in the USSR, German Democratic Republic, and China, often mimicking Western designs.
21st century
The 21st century witnessed the mass adoption of digital cameras and significant improvements in sensor technology. A major revolution came with the incorporation of cameras into smartphones, making photography a commonplace activity. The century also marked the rise of computational photography, using algorithms and AI to enhance image quality. Features like low-light and HDR photography, optical image stabilization, and depth-sensing became common in smartphone cameras.
Mechanics
Most cameras capture light from the visible spectrum, while specialized cameras capture other portions of the electromagnetic spectrum, such as infrared.
All cameras use the same basic design: light enters an enclosed box through a converging or convex lens and an image is recorded on a light-sensitive medium. A shutter mechanism controls the length of time that light enters the camera.
Most cameras also have a viewfinder, which shows the scene to be recorded, along with means to adjust various combinations of focus, aperture and shutter speed.
Exposure control
Aperture
Light enters the camera through an aperture, an opening adjusted by overlapping plates called the aperture ring. Typically located in the lens, this opening can be widened or narrowed to alter the amount of light that strikes the film or sensor. The size of the aperture can be set manually, by rotating the lens or adjusting a dial or automatically based on readings from an internal light meter.
As the aperture is adjusted, the opening expands and contracts in increments called f-stops. The smaller the f-stop, the more light is allowed to enter the lens, increasing the exposure. Typically, f-stops range from 1.4 to 32 in standard increments: 1.4, 2, 2.8, 4, 5.6, 8, 11, 16, 22, and 32. The light entering the camera is halved with each increasing increment.
The wider opening at lower f-stops narrows the range of focus so the background is blurry while the foreground is in focus. This depth of field increases as the aperture closes. A narrow aperture results in a high depth of field, meaning that objects at many different distances from the camera will appear to be in focus. What is acceptably in focus is determined by the circle of confusion, the photographic technique, the equipment in use and the degree of magnification expected of the final image.
Shutter
The shutter, along with the aperture, is one of two ways to control the amount of light entering the camera. The shutter determines the duration that the light-sensitive surface is exposed to light. The shutter opens, light enters the camera and exposes the film or sensor to light, and then the shutter closes.
There are two types of mechanical shutters: the leaf-type shutter and the focal-plane shutter. The leaf-type uses a circular iris diaphragm maintained under spring tension inside or just behind the lens that rapidly opens and closes when the shutter is released.
More commonly, a focal-plane shutter is used. This shutter operates close to the film plane and employs metal plates or cloth curtains with an opening that passes across the light-sensitive surface. The curtains or plates have an opening that is pulled across the film plane during exposure. The focal-plane shutter is typically used in single-lens reflex (SLR) cameras, since covering the film (rather than blocking the light passing through the lens) allows the photographer to view the image through the lens at all times, except during the exposure itself. Covering the film also facilitates removing the lens from a loaded camera, as many SLRs have interchangeable lenses.
A digital camera may use a mechanical or electronic shutter, the latter of which is common in smartphone cameras. Electronic shutters either record data from the entire sensor simultaneously (a global shutter) or record the data line by line across the sensor (a rolling shutter). In movie cameras, a rotary shutter opens and closes in sync with the advancement of each frame of film.
The duration for which the shutter is open is called the shutter speed or exposure time. Typical exposure times can range from one second to 1/1,000 of a second, though longer and shorter durations are not uncommon. In the early stages of photography, exposures were often several minutes long. These long exposure times often resulted in blurry images, as a single object is recorded in multiple places across a single image for the duration of the exposure. To prevent this, shorter exposure times can be used. Very short exposure times can capture fast-moving action and eliminate motion blur. However, shorter exposure times require more light to produce a properly exposed image, so shortening the exposure time is not always possible.
Like aperture settings, exposure times increment in powers of two. The two settings determine the exposure value (EV), a measure of how much light is recorded during the exposure. There is a direct relationship between the exposure times and aperture settings so that if the exposure time is lengthened one step, but the aperture opening is also narrowed one step, then the amount of light that contacts the film or sensor is the same.
Light meter
In most modern cameras, the amount of light entering the camera is measured using a built-in light meter or exposure meter. Taken through the lens (called metering), these readings are taken using a panel of light-sensitive semiconductors. They are used to calculate optimal exposure settings. These settings are typically determined automatically as the reading is used by the camera's microprocessor. The reading from the light meter is incorporated with aperture settings, exposure times, and film or sensor sensitivity to calculate the optimal exposure.
Light meters typically average the light in a scene to 18% middle gray. More advanced cameras are more nuanced in their metering—weighing the center of the frame more heavily (center-weighted metering), considering the differences in light across the image (matrix metering), or allowing the photographer to take a light reading at a specific point within the image (spot metering).
Lens
A camera lens is an assembly of multiple optical elements, typically made from high-quality glass. Its primary function is to focus light onto a camera's film or digital sensor, thereby producing an image. This process significantly influences image quality, the overall appearance of the photo, and which parts of the scene are brought into focus.
A camera lens is constructed from a series of lens elements, small pieces of glass arranged to form an image accurately on the light-sensitive surface. Each element is designed to reduce optical aberrations, or distortions, such as chromatic aberration (a failure of the lens to focus all colors at the same point), vignetting (darkening of image corners), and distortion (bending or warping of the image). The degree of these distortions can vary depending on the subject of the photo.
The focal length of the lens, measured in millimeters, plays a critical role as it determines how much of the scene the camera can capture and how large the objects appear. Wide-angle lenses provide a broad view of the scene, while telephoto lenses capture a narrower view but magnify the objects. The focal length also influences the ease of taking clear pictures handheld, with longer lengths making it more challenging to avoid blur from small camera movements.
Two primary types of lenses include zoom and prime lenses. A zoom lens allows for changing its focal length within a certain range, providing the convenience of adjusting the scene capture without moving the camera or changing the lens. A prime lens, in contrast, has a fixed focal length. While less flexible, prime lenses often provide superior image quality, are typically lighter, and perform better in low light.
Focus involves adjusting the lens elements to sharpen the image of the subject at various distances. The focus is adjusted through the focus ring on the lens, which moves the lens elements closer or further from the sensor. Autofocus is a feature included in many lenses, which uses a motor within the lens to adjust the focus quickly and precisely based on the lens's detection of contrast or phase differences. This feature can be enabled or disabled using switches on the lens body.
Advanced lenses may include mechanical image stabilization systems that move lens elements or the image sensor itself to counteract camera shake, especially beneficial in low-light conditions or at slow shutter speeds. Lens hoods, filters, and caps are accessories used alongside a lens to enhance image quality, protect the lens, or achieve specific effects.
Viewfinder
The camera's viewfinder provides a real-time approximation of what will be captured by the sensor or film. It assists photographers in aligning, focusing, and adjusting the composition, lighting, and exposure of their shots, enhancing the accuracy of the final image.
Viewfinders fall into two primary categories: optical and electronic. Optical viewfinders, commonly found in Single-Lens Reflex (SLR) cameras, use a system of mirrors or prisms to reflect light from the lens to the viewfinder, providing a clear, real-time view of the scene. Electronic viewfinders, typical in mirrorless cameras, project an electronic image onto a small display, offering a wider range of information such as live exposure previews and histograms, albeit at the cost of potential lag and higher battery consumption. Specialized viewfinder systems exist for specific applications, like subminiature cameras for spying or underwater photography.
Parallax error, resulting from misalignment between the viewfinder and lens axes, can cause inaccurate representations of the subject's position. While negligible with distant subjects, this error becomes prominent with closer ones. Some viewfinders incorporate parallax-compensating devices to mitigate that issue.
Film and sensor
Image capture in a camera occurs when light strikes a light-sensitive surface: photographic film or a digital sensor. Housed within the camera body, the film or sensor records the light's pattern when the shutter is briefly opened to allow light to pass during the exposure.
Loading film into a film camera is a manual process. The film, typically housed in a cartridge, is loaded into a designated slot in the camera. One end of the film strip, the film leader, is manually threaded onto a take-up spool. Once the back of the camera is closed, the film advance lever or knob is used to ensure the film is correctly placed. The photographer then winds the film, either manually or automatically depending on the camera, to position a blank portion of the film in the path of the light. Each time a photo is taken, the film advance mechanism moves the exposed film out of the way, bringing a new, unexposed section of film into position for the next shot.
The film must be advanced after each shot to prevent double exposure — where the same section of film is exposed to light twice, resulting in overlapped images. Once all frames on the film roll have been exposed, the film is rewound back into the cartridge, ready to be removed from the camera for developing.
In digital cameras, sensors typically comprise Charge-Coupled Devices (CCDs) or Complementary Metal-Oxide-Semiconductor (CMOS) chips, both of which convert incoming light into electrical charges to form digital images. CCD sensors, though power-intensive, are recognized for their excellent light sensitivity and image quality. Conversely, CMOS sensors offer individual pixel readouts, leading to less power consumption and faster frame rates, with their image quality having improved significantly over time.
Digital cameras convert light into electronic data that can be directly processed and stored. The volume of data generated is dictated by the sensor's size and properties, necessitating storage media such as Compact Flash, Memory Sticks, and SD (Secure Digital) cards. Modern digital cameras typically feature a built-in monitor for immediate image review and adjustments. Digital images are also more readily handled and manipulated by computers, offering a significant advantage in terms of flexibility and post-processing potential over traditional film.
Camera accessories
Flash
A flash provides a short burst of bright light during exposure and is a commonly used artificial light source in photography. Most modern flash systems use a battery-powered high-voltage discharge through a gas-filled tube to generate bright light for a very short time (1/1,000 of a second or less).
Many flash units measure the light reflected from the flash to help determine the appropriate duration of the flash. When the flash is attached directly to the camera—typically in a slot at the top of the camera (the flash shoe or hot shoe) or through a cable—activating the shutter on the camera triggers the flash, and the camera's internal light meter can help determine the duration of the flash.
Additional flash equipment can include a light diffuser, mount and stand, reflector, soft box, trigger and cord.
Other accessories
Accessories for cameras are mainly used for care, protection, special effects, and functions.
Lens hood: used on the end of a lens to block the sun or other light source to prevent glare and lens flare (see also matte box).
Lens cap: covers and protects the camera lens when not in use.
Lens adapter: allows the use of lenses other than those for which the camera was designed.
Filter: allows artificial colors or changes light density.
Lens extension tube: allows close focus in macro photography.
Care and protection: including camera case and cover, maintenance tools, and screen protector.
Camera monitor: provides an off-camera view of the composition with a brighter and more colorful screen, and typically exposes more advanced tools such as framing guides, focus peaking, zebra stripes, waveform monitors (oftentimes as an "RGB parade"), vectorscopes and false color to highlight areas of the image critical to the photographer.
Tripod: primarily used for keeping the camera steady while recording video, doing a long exposure, and time-lapse photography.
Microscope adapter: used to connect a camera to a microscope to photograph what the microscope is examining.
Cable release: used to remotely control the shutter using a remote shutter button that can be connected to the camera via a cable. It can be used to lock the shutter open for the desired period, and it is also commonly used to prevent the camera shake from pressing the built-in camera shutter button.
Dew shield: prevents moisture build-up on the lens.
UV filter: can protect the front element of a lens from scratches, cracks, smudges, dirt, dust, and moisture while keeping a minimum impact on image quality.
Battery and sometimes a charger.
Large format cameras use special equipment that includes a magnifier loupe, view finder, angle finder, and focusing rail/truck. Some professional SLRs can be provided with interchangeable finders for eye-level or waist-level focusing, focusing screens, eyecup, data backs, motor-drives for film transportation or external battery packs.
Primary types
Single-lens reflex (SLR) camera
In photography, the single-lens reflex camera (SLR) is provided with a mirror to redirect light from the lens to the viewfinder prior to releasing the shutter for composing and focusing an image. When the shutter is released, the mirror swings up and away, allowing the exposure of the photographic medium, and instantly returns after the exposure is finished. No SLR camera before 1954 had this feature, although the mirror on some early SLR cameras was entirely operated by the force exerted on the shutter release and only returned when the finger pressure was released. The Asahiflex II, released by Japanese company Asahi (Pentax) in 1954, was the world's first SLR camera with an instant return mirror.
In the single-lens reflex camera, the photographer sees the scene through the camera lens. This avoids the problem of parallax which occurs when the viewfinder or viewing lens is separated from the taking lens. Single-lens reflex cameras have been made in several formats including sheet film 5x7" and 4x5", roll film 220/120 taking 8,10, 12, or 16 photographs on a 120 roll, and twice that number of a 220 film. These correspond to 6x9, 6x7, 6x6, and 6x4.5 respectively (all dimensions in cm). Notable manufacturers of large format and roll film SLR cameras include Bronica, Graflex, Hasselblad, Seagull, Mamiya and Pentax. However, the most common format of SLR cameras has been 35 mm and subsequently the migration to digital SLR cameras, using almost identical sized bodies and sometimes using the same lens systems.
Almost all SLR cameras use a front-surfaced mirror in the optical path to direct the light from the lens via a viewing screen and pentaprism to the eyepiece. At the time of exposure, the mirror is flipped up out of the light path before the shutter opens. Some early cameras experimented with other methods of providing through-the-lens viewing, including the use of a semi-transparent pellicle as in the Canon Pellix and others with a small periscope such as in the Corfield Periflex series.
Large-format camera
The large-format camera, taking sheet film, is a direct successor of the early plate cameras and remained in use for high-quality photography and technical, architectural, and industrial photography. There are three common types: the view camera, with its monorail and field camera variants, and the press camera. They have extensible bellows with the lens and shutter mounted on a lens plate at the front. Backs taking roll film and later digital backs are available in addition to the standard dark slide back. These cameras have a wide range of movements allowing very close control of focus and perspective. Composition and focusing are done on view cameras by viewing a ground-glass screen which is replaced by the film to make the exposure; they are suitable for static subjects only and are slow to use.
Plate camera
The earliest cameras produced in significant numbers were plate cameras, using sensitized glass plates. Light entered a lens mounted on a lens board which was separated from the plate by extendible bellows. There were simple box cameras for glass plates but also single-lens reflex cameras with interchangeable lenses and even for color photography (Autochrome Lumière). Many of these cameras had controls to raise, lower, and tilt the lens forwards or backward to control perspective.
Focusing of these plate cameras was by the use of a ground glass screen at the point of focus. Because lens design only allowed rather small aperture lenses, the image on the ground glass screen was faint and most photographers had a dark cloth to cover their heads to allow focusing and composition to be carried out more quickly. When focus and composition were satisfactory, the ground glass screen was removed, and a sensitized plate was put in its place protected by a dark slide. To make the exposure, the dark decline was carefully slid out and the shutter opened, and then closed and the dark fall replaced.
Glass plates were later replaced by sheet film in a dark slide for sheet film; adapter sleeves were made to allow sheet film to be used in plate holders. In addition to the ground glass, a simple optical viewfinder was often fitted.
Medium-format camera
Medium-format cameras have a film size between the large-format cameras and smaller 35 mm cameras. Typically these systems use 120 or 220 roll film. The most common image sizes are 6×4.5 cm, 6×6 cm and 6×7 cm; the older 6×9 cm is rarely used. The designs of this kind of camera show greater variation than their larger brethren, ranging from monorail systems through the classic Hasselblad model with separate backs, to smaller rangefinder cameras. There are even compact amateur cameras available in this format.
Twin-lens reflex camera
Twin-lens reflex cameras used a pair of nearly identical lenses: one to form the image and one as a viewfinder. The lenses were arranged with the viewing lens immediately above the taking lens. The viewing lens projects an image onto a viewing screen which can be seen from above. Some manufacturers such as Mamiya also provided a reflex head to attach to the viewing screen to allow the camera to be held to the eye when in use. The advantage of a TLR was that it could be easily focused using the viewing screen and that under most circumstances the view seen on the viewing screen was identical to that recorded on film. At close distances, however, parallax errors were encountered, and some cameras also included an indicator to show what part of the composition would be excluded.
Some TLRs had interchangeable lenses, but as these had to be paired lenses, they were relatively heavy and did not provide the range of focal lengths that the SLR could support. Most TLRs used 120 or 220 films; some used the smaller 127 films.
Compact cameras
Instant camera
After exposure, every photograph is taken through pinch rollers inside the instant camera. Thereby the developer paste contained in the paper 'sandwich' is distributed on the image. After a minute, the cover sheet just needs to be removed and one gets a single original positive image with a fixed format. With some systems, it was also possible to create an instant image negative, from which then could be made copies in the photo lab. The ultimate development was the SX-70 system of Polaroid, in which a row of ten shots – engine driven – could be made without having to remove any cover sheets from the picture. There were instant cameras for a variety of formats, as well as adapters for instant film use in medium- and large-format cameras.
Subminiature camera
Subminiature cameras were first produced in the twentieth century and use film significantly smaller than 35mm. The expensive 8×11mm Minox, the only type of camera produced by the company from 1937 to 1976, became very widely known and was often used for espionage (the Minox company later also produced larger cameras). Later inexpensive subminiatures were made for general use, some using rewound 16 mm cine film. Image quality with these small film sizes was limited.
Folding camera
The introduction of films enabled the existing designs for plate cameras to be made much smaller and for the baseplate to be hinged so that it could be folded up, compressing the bellows. These designs were very compact and small models were dubbed vest pocket cameras. One of the smallest and best-selling cameras was the Vest Pocket Kodak, sold in two generations between 1912 and 1934. Folding roll film cameras were preceded by folding plate cameras, more compact than other designs.
Box camera
Box cameras were introduced as budget-level cameras and had few if any controls. The original box Brownie models had a small reflex viewfinder mounted on the top of the camera and had no aperture or focusing controls and just a simple shutter. Later models such as the Brownie 127 had larger direct view optical viewfinders together with a curved film path to reduce the impact of deficiencies in the lens.
Rangefinder camera
As camera lens technology developed and wide aperture lenses became more common, rangefinder cameras were introduced to make focusing more precise. Early rangefinders had two separate viewfinder windows, one of which is linked to the focusing mechanisms and moved right or left as the focusing ring is turned. The two separate images are brought together on a ground glass viewing screen. When vertical lines in the object being photographed meet exactly in the combined image, the object is in focus. A normal composition viewfinder is also provided. Later the viewfinder and rangefinder were combined. Many rangefinder cameras had interchangeable lenses, each lens requiring its range- and viewfinder linkages.
Rangefinder cameras were produced in half- and full-frame 35 mm and roll film (medium format).
Motion picture cameras
A movie camera or a video camera operates similarly to a still camera, except it records a series of static images in rapid succession, commonly at a rate of 24 frames per second. When the images are combined and displayed in order, the illusion of motion is achieved.
Cameras that capture many images in sequence are known as movie cameras or as cine cameras in Europe; those designed for single images are still cameras. However, these categories overlap as still cameras are often used to capture moving images in special effects work and many modern cameras can quickly switch between still and motion recording modes.
A ciné camera or movie camera takes a rapid sequence of photographs on an image sensor or strips of film. In contrast to a still camera, which captures a single snapshot at a time, the ciné camera takes a series of images, each called a frame, through the use of an intermittent mechanism.
The frames are later played back in a ciné projector at a specific speed, called the frame rate (number of frames per second). While viewing, a person's visual system merge the separate pictures to create the illusion of motion. The first ciné camera was built around 1888 and by 1890 several types were being manufactured. The standard film size for ciné cameras was quickly established as 35mm film and this remained in use until the transition to digital cinematography. Other professional standard formats include 70 mm film and 16 mm film whilst amateur filmmakers used 9.5 mm film, 8 mm film, or Standard 8 and Super 8 before the move into digital format.
The size and complexity of ciné cameras vary greatly depending on the uses required of the camera. Some professional equipment is very large and too heavy to be handheld whilst some amateur cameras were designed to be very small and light for single-handed operation.
Professional video camera
A professional video camera (often called a television camera even though the use has spread beyond television) is a high-end device for creating electronic moving images (as opposed to a movie camera, that earlier recorded the images on film). Originally developed for use in television studios, they are now also used for music videos, direct-to-video movies, corporate and educational videos, marriage videos, etc.
These cameras earlier used vacuum tubes and later electronic image sensors.
Camcorders
A camcorder is an electronic device combining a video camera and a video recorder. Although marketing materials may use the colloquial term "camcorder", the name on the package and manual is often "video camera recorder". Most devices capable of recording video are camera phones and digital cameras primarily intended for still pictures; the term "camcorder" is used to describe a portable, self-contained device, with video capture and recording its primary function.
Digital camera
A digital camera (or digicam) is a camera that encodes digital images and videos and stores them for later reproduction. They typically use semiconductor image sensors. Most cameras sold today are digital, and they are incorporated into many devices ranging from mobile phones (called camera phones) to vehicles.
Digital and film cameras share an optical system, typically using a lens of variable aperture to focus light onto an image pickup device. The aperture and shutter admit the correct amount of light to the imager, just as with film but the image pickup device is electronic rather than chemical. However, unlike film cameras, digital cameras can display images on a screen immediately after being captured or recorded, and store and delete images from memory. Most digital cameras can also record moving videos with sound. Some digital cameras can crop and stitch pictures & perform other elementary image editing.
Consumers adopted digital cameras in the 1990s. Professional video cameras transitioned to digital around the 2000s–2010s. Finally, movie cameras transitioned to digital in the 2010s.
The first camera using digital electronics to capture and store images was developed by Kodak engineer Steven Sasson in 1975. He used a charge-coupled device (CCD) provided by Fairchild Semiconductor, which provided only 0.01 megapixels to capture images. Sasson combined the CCD device with movie camera parts to create a digital camera that saved black and white images onto a cassette tape.The images were then read from the cassette and viewed on a TV monitor. Later, cassette tapes were replaced by flash memory.
In 1986, Japanese company Nikon introduced an analog-recording electronic single-lens reflex camera, the Nikon SVC.
The first full-frame digital SLR cameras were developed in Japan from around 2000 to 2002: the MZ-D by Pentax, the N Digital by Contax's Japanese R6D team, and the EOS-1Ds by Canon. Gradually in the 2000s, the full-frame DSLR became the dominant camera type for professional photography.
On most digital cameras a display, often a liquid crystal display (LCD), permits the user to view the scene to be recorded and settings such as ISO speed, exposure, and shutter speed.
Camera phone
In 2000, Sharp introduced the world's first digital camera phone, the J-SH04 J-Phone, in Japan. By the mid-2000s, higher-end cell phones had an integrated digital camera, and by the beginning of the 2010s, almost all smartphones had an integrated digital camera.
| Technology | Optical | null |
52649 | https://en.wikipedia.org/wiki/Acetylcholine | Acetylcholine | Acetylcholine (ACh) is an organic compound that functions in the brain and body of many types of animals (including humans) as a neurotransmitter. Its name is derived from its chemical structure: it is an ester of acetic acid and choline. Parts in the body that use or are affected by acetylcholine are referred to as cholinergic.
Acetylcholine is the neurotransmitter used at the neuromuscular junction—in other words, it is the chemical that motor neurons of the nervous system release in order to activate muscles. This property means that drugs that affect cholinergic systems can have very dangerous effects ranging from paralysis to convulsions. Acetylcholine is also a neurotransmitter in the autonomic nervous system, both as an internal transmitter for both the sympathetic and the parasympathetic nervous system, and as the final product released by the parasympathetic nervous system. Acetylcholine is the primary neurotransmitter of the parasympathetic nervous system.
In the brain, acetylcholine functions as a neurotransmitter and as a neuromodulator. The brain contains a number of cholinergic areas, each with distinct functions; such as playing an important role in arousal, attention, memory and motivation. Acetylcholine has also been found in cells of non-neural origins as well as microbes. Recently, enzymes related to its synthesis, degradation and cellular uptake have been traced back to early origins of unicellular eukaryotes. The protist pathogens Acanthamoeba spp. have shown evidence of the presence of ACh, which provides growth and proliferative signals via a membrane-located M1-muscarinic receptor homolog.
Partly because of acetylcholine's muscle-activating function, but also because of its functions in the autonomic nervous system and brain, many important drugs exert their effects by altering cholinergic transmission. Numerous venoms and toxins produced by plants, animals, and bacteria, as well as chemical nerve agents such as sarin, cause harm by inactivating or hyperactivating muscles through their influences on the neuromuscular junction. Drugs that act on muscarinic acetylcholine receptors, such as atropine, can be poisonous in large quantities, but in smaller doses they are commonly used to treat certain heart conditions and eye problems. Scopolamine, or diphenhydramine, which also act mainly on muscarinic receptors in an inhibitory fashion in the brain (especially the M1 receptor) can cause delirium, hallucinations, and amnesia through receptor antagonism at these sites. So far as of 2016, only the M1 receptor subtype has been implicated in anticholinergic delirium. The addictive qualities of nicotine are derived from its effects on nicotinic acetylcholine receptors in the brain.
Chemistry
Acetylcholine is a choline molecule that has been acetylated at the oxygen atom. Because of the charged ammonium group, acetylcholine does not penetrate lipid membranes. Because of this, when the molecule is introduced externally, it remains in the extracellular space and at present it is considered that the molecule does not pass through the blood–brain barrier.
Biochemistry
Acetylcholine is synthesized in certain neurons by the enzyme choline acetyltransferase from the compounds choline and acetyl-CoA. Cholinergic neurons are capable of producing ACh. An example of a central cholinergic area is the nucleus basalis of Meynert in the basal forebrain.
The enzyme acetylcholinesterase converts acetylcholine into the inactive metabolites choline and acetate. This enzyme is abundant in the synaptic cleft, and its role in rapidly clearing free acetylcholine from the synapse is essential for proper muscle function. Certain neurotoxins work by inhibiting acetylcholinesterase, thus leading to excess acetylcholine at the neuromuscular junction, causing paralysis of the muscles needed for breathing and stopping the beating of the heart.
Functions
Acetylcholine functions in both the central nervous system (CNS) and the peripheral nervous system (PNS). In the CNS, cholinergic projections from the basal forebrain to the cerebral cortex and hippocampus support the cognitive functions of those target areas. In the PNS, acetylcholine activates muscles and is a major neurotransmitter in the autonomic nervous system.
Cellular effects
Like many other biologically active substances, acetylcholine exerts its effects by binding to and activating receptors located on the surface of cells. There are two main classes of acetylcholine receptor, nicotinic and muscarinic. They are named for chemicals that can selectively activate each type of receptor without activating the other: muscarine is a compound found in the mushroom Amanita muscaria; nicotine is found in tobacco.
Nicotinic acetylcholine receptors are ligand-gated ion channels permeable to sodium, potassium, and calcium ions. In other words, they are ion channels embedded in cell membranes, capable of switching from a closed to an open state when acetylcholine binds to them; in the open state they allow ions to pass through. Nicotinic receptors come in two main types, known as muscle-type and neuronal-type. The muscle-type can be selectively blocked by curare, the neuronal-type by hexamethonium. The main location of muscle-type receptors is on muscle cells, as described in more detail below. Neuronal-type receptors are located in autonomic ganglia (both sympathetic and parasympathetic), and in the central nervous system.
Muscarinic acetylcholine receptors have a more complex mechanism, and affect target cells over a longer time frame. In mammals, five subtypes of muscarinic receptors have been identified, labeled M1 through M5. All of them function as G protein-coupled receptors, meaning that they exert their effects via a second messenger system. The M1, M3, and M5 subtypes are Gq-coupled; they increase intracellular levels of IP3 and calcium by activating phospholipase C. Their effect on target cells is usually excitatory. The M2 and M4 subtypes are Gi/Go-coupled; they decrease intracellular levels of cAMP by inhibiting adenylate cyclase. Their effect on target cells is usually inhibitory. Muscarinic acetylcholine receptors are found in both the central nervous system and the peripheral nervous system of the heart, lungs, upper gastrointestinal tract, and sweat glands.
Neuromuscular junction
Acetylcholine is the substance the nervous system uses to activate skeletal muscles, a kind of striated muscle. These are the muscles used for all types of voluntary movement, in contrast to smooth muscle tissue, which is involved in a range of involuntary activities such as movement of food through the gastrointestinal tract and constriction of blood vessels. Skeletal muscles are directly controlled by motor neurons located in the spinal cord or, in a few cases, the brainstem. These motor neurons send their axons through motor nerves, from which they emerge to connect to muscle fibers at a special type of synapse called the neuromuscular junction.
When a motor neuron generates an action potential, it travels rapidly along the nerve until it reaches the neuromuscular junction, where it initiates an electrochemical process that causes acetylcholine to be released into the space between the presynaptic terminal and the muscle fiber. The acetylcholine molecules then bind to nicotinic ion-channel receptors on the muscle cell membrane, causing the ion channels to open. Sodium ions then flow into the muscle cell, initiating a sequence of steps that finally produce muscle contraction.
Factors that decrease release of acetylcholine (and thereby affecting P-type calcium channels):
Antibiotics (clindamycin, polymyxin)
Magnesium: antagonizes P-type calcium channels
Hypocalcemia
Anticonvulsants
Diuretics (furosemide)
Eaton-Lambert syndrome: inhibits P-type calcium channels
Myasthenia gravis
Botulinum toxin: inhibits SNARE proteins
Calcium channel blockers (nifedipine, diltiazem) do not affect P-channels. These drugs affect L-type calcium channels.
Autonomic nervous system
The autonomic nervous system controls a wide range of involuntary and unconscious body functions. Its main branches are the sympathetic nervous system and parasympathetic nervous system. Broadly speaking, the function of the sympathetic nervous system is to mobilize the body for action; the phrase often invoked to describe it is fight-or-flight. The function of the parasympathetic nervous system is to put the body in a state conducive to rest, regeneration, digestion, and reproduction; the phrase often invoked to describe it is "rest and digest" or "feed and breed". Both of these aforementioned systems use acetylcholine, but in different ways.
At a schematic level, the sympathetic and parasympathetic nervous systems are both organized in essentially the same way: preganglionic neurons in the central nervous system send projections to neurons located in autonomic ganglia, which send output projections to virtually every tissue of the body. In both branches the internal connections, the projections from the central nervous system to the autonomic ganglia, use acetylcholine as a neurotransmitter to innervate (or excite) ganglia neurons. In the parasympathetic nervous system the output connections, the projections from ganglion neurons to tissues that do not belong to the nervous system, also release acetylcholine but act on muscarinic receptors. In the sympathetic nervous system the output connections mainly release noradrenaline, although acetylcholine is released at a few points, such as the sudomotor innervation of the sweat glands.
Direct vascular effects
Acetylcholine in the serum exerts a direct effect on vascular tone by binding to muscarinic receptors present on vascular endothelium. These cells respond by increasing production of nitric oxide, which signals the surrounding smooth muscle to relax, leading to vasodilation.
Central nervous system
In the central nervous system, ACh has a variety of effects on plasticity, arousal and reward. ACh has an important role in the enhancement of alertness when we wake up, in sustaining attention and in learning and memory.
Damage to the cholinergic (acetylcholine-producing) system in the brain has been shown to be associated with the memory deficits associated with Alzheimer's disease. ACh has also been shown to promote REM sleep.
In the brainstem acetylcholine originates from the Pedunculopontine nucleus and laterodorsal tegmental nucleus collectively known as the mesopontine tegmentum area or pontomesencephalotegmental complex. In the basal forebrain, it originates from the basal nucleus of Meynert and medial septal nucleus:
The pontomesencephalotegmental complex acts mainly on M1 receptors in the brainstem, deep cerebellar nuclei, pontine nuclei, locus coeruleus, raphe nucleus, lateral reticular nucleus and inferior olive. It also projects to the thalamus, tectum, basal ganglia and basal forebrain.
Basal nucleus of Meynert acts mainly on M1 receptors in the neocortex.
Medial septal nucleus acts mainly on M1 receptors in the hippocampus and parts of the cerebral cortex.
In addition, ACh acts as an important internal transmitter in the striatum, which is part of the basal ganglia. It is released by cholinergic interneurons. In humans, non-human primates and rodents, these interneurons respond to salient environmental stimuli with responses that are temporally aligned with the responses of dopaminergic neurons of the substantia nigra.
Memory
Acetylcholine has been implicated in learning and memory in several ways. The anticholinergic drug scopolamine impairs acquisition of new information in humans and animals. In animals, disruption of the supply of acetylcholine to the neocortex impairs the learning of simple discrimination tasks, comparable to the acquisition of factual information and disruption of the supply of acetylcholine to the hippocampus and adjacent cortical areas produces forgetfulness, comparable to anterograde amnesia in humans.
Diseases and disorders
Myasthenia gravis
The disease myasthenia gravis, characterized by muscle weakness and fatigue, occurs when the body inappropriately produces antibodies against acetylcholine nicotinic receptors, and thus inhibits proper acetylcholine signal transmission. Over time, the motor end plate is destroyed. Drugs that competitively inhibit acetylcholinesterase (e.g., neostigmine, physostigmine, or primarily pyridostigmine) are effective in treating the symptoms of this disorder. They allow endogenously released acetylcholine more time to interact with its respective receptor before being inactivated by acetylcholinesterase in the synaptic cleft (the space between nerve and muscle).
Pharmacology
Blocking, hindering or mimicking the action of acetylcholine has many uses in medicine. Drugs acting on the acetylcholine system are either agonists to the receptors, stimulating the system, or antagonists, inhibiting it. Acetylcholine receptor agonists and antagonists can either have an effect directly on the receptors or exert their effects indirectly, e.g., by affecting the enzyme acetylcholinesterase, which degrades the receptor ligand. Agonists increase the level of receptor activation; antagonists reduce it.
Acetylcholine itself does not have therapeutic value as a drug for intravenous administration because of its multi-faceted action (non-selective) and rapid inactivation by cholinesterase. However, it is used in the form of eye drops to cause constriction of the pupil during cataract surgery, which facilitates quick post-operational recovery.
Nicotinic receptors
Nicotine binds to and activates nicotinic acetylcholine receptors, mimicking the effect of acetylcholine at these receptors. ACh opens a Na+ channel upon binding so that Na+ flows into the cell. This causes a depolarization, and results in an excitatory post-synaptic potential. Thus, ACh is excitatory on skeletal muscle; the electrical response is fast and short-lived. Curares are arrow poisons, which act at nicotinic receptors and have been used to develop clinically useful therapies.
Muscarinic receptors
Muscarinic receptors form G protein-coupled receptor complexes in the cell membranes of neurons and other cells. Atropine is a non-selective competitive antagonist with Acetylcholine at muscarinic receptors.
Cholinesterase inhibitors
Many ACh receptor agonists work indirectly by inhibiting the enzyme acetylcholinesterase. The resulting accumulation of acetylcholine causes continuous stimulation of the muscles, glands, and central nervous system, which can result in fatal convulsions if the dose is high.
They are examples of enzyme inhibitors, and increase the action of acetylcholine by delaying its degradation; some have been used as nerve agents (Sarin and VX nerve gas) or pesticides (organophosphates and the carbamates). Many toxins and venoms produced by plants and animals also contain cholinesterase inhibitors. In clinical use, they are administered in low doses to reverse the action of muscle relaxants, to treat myasthenia gravis, and to treat symptoms of Alzheimer's disease (rivastigmine, which increases cholinergic activity in the brain).
Synthesis inhibitors
Organic mercurial compounds, such as methylmercury, have a high affinity for sulfhydryl groups, which causes dysfunction of the enzyme choline acetyltransferase. This inhibition may lead to acetylcholine deficiency, and can have consequences on motor function.
Release inhibitors
Botulinum toxin (Botox) acts by suppressing the release of acetylcholine, whereas the venom from a black widow spider (alpha-latrotoxin) has the reverse effect. ACh inhibition causes paralysis. When bitten by a black widow spider, one experiences the wastage of ACh supplies and the muscles begin to contract. If and when the supply is depleted, paralysis occurs.
Comparative biology and evolution
Acetylcholine is used by organisms in all domains of life for a variety of purposes. It is believed that choline, a precursor to acetylcholine, was used by single celled organisms billions of years ago for synthesizing cell membrane phospholipids. Following the evolution of choline transporters, the abundance of intracellular choline paved the way for choline to become incorporated into other synthetic pathways, including acetylcholine production. Acetylcholine is used by bacteria, fungi, and a variety of other animals. Many of the uses of acetylcholine rely on its action on ion channels via GPCRs like membrane proteins.
The two major types of acetylcholine receptors, muscarinic and nicotinic receptors, have convergently evolved to be responsive to acetylcholine. This means that rather than having evolved from a common homolog, these receptors evolved from separate receptor families. It is estimated that the nicotinic receptor family dates back longer than 2.5 billion years. Likewise, muscarinic receptors are thought to have diverged from other GPCRs at least 0.5 billion years ago. Both of these receptor groups have evolved numerous subtypes with unique ligand affinities and signaling mechanisms. The diversity of the receptor types enables acetylcholine to create varying responses depending on which receptor types are activated, and allow for acetylcholine to dynamically regulate physiological processes. ACh receptors are related to 5-HT3 (serotonin), GABA, and Glycine receptors, both in sequence and structure, strongly suggesting that they have a common evolutionary origin.
History
In 1867, Adolf von Baeyer resolved the structures of choline and acetylcholine and synthesized them both, referring to the latter as acetylneurin in the study. Choline is a precursor for acetylcholine. Acetylcholine was first noted to be biologically active in 1906, when Reid Hunt (1870–1948) and René de M. Taveau found that it decreased blood pressure in exceptionally tiny doses. This was after Frederick Walker Mott and William Dobinson Halliburton noted in 1899 that choline injections decreased the blood pressure of animals.
In 1914, Arthur J. Ewins was the first to extract acetylcholine from nature. He identified it as the blood pressure-decreasing contaminant from some Claviceps purpurea ergot extracts, by the request of Henry Hallett Dale. Later in 1914, Dale outlined the effects of acetylcholine at various types of peripheral synapses and also noted that it lowered the blood pressure of cats via subcutaneous injections even at doses of one nanogram.
The concept of neurotransmitters was unknown until 1921, when Otto Loewi noted that the vagus nerve secreted a substance that inhibited the heart muscle whilst working as a professor in the University of Graz. He named it vagusstoff ("vagus substance"), noted it to be a structural analog of choline and suspected it to be acetylcholine. In 1926, Loewi and E. Navratil deduced that the compound is probably acetylcholine, as vagusstoff and synthetic acetylcholine lost their activity in a similar manner when in contact with tissue lysates that contained acetylcholine-degrading enzymes (now known to be cholinesterases). This conclusion was accepted widely. Later studies confirmed the function of acetylcholine as a neurotransmitter.
In 1936, H. H. Dale and O. Loewi shared the Nobel Prize in Physiology or Medicine for their studies of acetylcholine and nerve impulses.
| Biology and health sciences | Neurotransmitters | Biology |
52713 | https://en.wikipedia.org/wiki/Binary%20star | Binary star | A binary star or binary star system is a system of two stars that are gravitationally bound to and in orbit around each other. Binary stars in the night sky that are seen as a single object to the naked eye are often resolved as separate stars using a telescope, in which case they are called visual binaries. Many visual binaries have long orbital periods of several centuries or millennia and therefore have orbits which are uncertain or poorly known. They may also be detected by indirect techniques, such as spectroscopy (spectroscopic binaries) or astrometry (astrometric binaries). If a binary star happens to orbit in a plane along our line of sight, its components will eclipse and transit each other; these pairs are called eclipsing binaries, or, together with other binaries that change brightness as they orbit, photometric binaries.
If components in binary star systems are close enough, they can gravitationally distort each other's outer stellar atmospheres. In some cases, these close binary systems can exchange mass, which may bring their evolution to stages that single stars cannot attain. Examples of binaries are Sirius, and Cygnus X-1 (Cygnus X-1 being a well-known black hole). Binary stars are also common as the nuclei of many planetary nebulae, and are the progenitors of both novae and type Ia supernovae.
Discovery
Double stars, a pair of stars that appear close to each other, have been observed since the invention of the telescope. Early examples include Mizar and Acrux. Mizar, in the Big Dipper (Ursa Major), was observed to be double by Giovanni Battista Riccioli in 1650 (and probably earlier by Benedetto Castelli and Galileo). The bright southern star Acrux, in the Southern Cross, was discovered to be double by Father Fontenay in 1685.
Evidence that stars in pairs were more than just optical alignments came in 1767 when English natural philosopher and clergyman John Michell became the first person to apply the mathematics of statistics to the study of the stars, demonstrating in a paper that many more stars occur in pairs or groups than a perfectly random distribution and chance alignment could account for. He focused his investigation on the Pleiades cluster, and calculated that the likelihood of finding such a close grouping of stars was about one in half a million. He concluded that the stars in these double or multiple star systems might be drawn to one another by gravitational pull, thus providing the first evidence for the existence of binary stars and star clusters.
William Herschel began observing double stars in 1779, hoping to find a near star paired with a distant star so he could measure the near star's changing position as the Earth orbited the Sun (measure its parallax), allowing him to calculate the distance to the near star. He would soon publish catalogs of about 700 double stars. By 1803, he had observed changes in the relative positions in a number of double stars over the course of 25 years, and concluded that, instead of showing parallax changes, they seemed to be orbiting each other in binary systems. The first orbit of a binary star was computed in 1827, when Félix Savary computed the orbit of Xi Ursae Majoris.
Over the years, many more double stars have been catalogued and measured. As of June 2017, the Washington Double Star Catalog, a database of visual double stars compiled by the United States Naval Observatory, contains over 100,000 pairs of double stars, including optical doubles as well as binary stars. Orbits are known for only a few thousand of these double stars.
Etymology
The term binary was first used in this context by Sir William Herschel in 1802, when he wrote:
By the modern definition, the term binary star is generally restricted to pairs of stars which revolve around a common center of mass. Binary stars which can be resolved with a telescope or interferometric methods are known as visual binaries. For most of the known visual binary stars one whole revolution has not been observed yet; rather, they are observed to have travelled along a curved path or a partial arc.
The more general term double star is used for pairs of stars which are seen to be close together in the sky. This distinction is rarely made in languages other than English. Double stars may be binary systems or may be merely two stars that appear to be close together in the sky but have vastly different true distances from the Sun. The latter are termed optical doubles or optical pairs.
Classifications
Methods of observation
Binary stars are classified into four types according to the way in which they are observed: visually, by observation; spectroscopically, by periodic changes in spectral lines; photometrically, by changes in brightness caused by an eclipse; or astrometrically, by measuring a deviation in a star's position caused by an unseen companion. Any binary star can belong to several of these classes; for example, several spectroscopic binaries are also eclipsing binaries.
Visual binaries
A visual binary star is a binary star for which the angular separation between the two components is great enough to permit them to be observed as a double star in a telescope, or even high-powered binoculars. The angular resolution of the telescope is an important factor in the detection of visual binaries, and as better angular resolutions are applied to binary star observations, an increasing number of visual binaries will be detected. The relative brightness of the two stars is also an important factor, as glare from a bright star may make it difficult to detect the presence of a fainter component.
The brighter star of a visual binary is the primary star, and the dimmer is considered the secondary. In some publications (especially older ones), a faint secondary is called the comes (plural comites; companion). If the stars are the same brightness, the discoverer designation for the primary is customarily accepted.
The position angle of the secondary with respect to the primary is measured, together with the angular distance between the two stars. The time of observation is also recorded. After a sufficient number of observations are recorded over a period of time, they are plotted in polar coordinates with the primary star at the origin, and the most probable ellipse is drawn through these points such that the Keplerian law of areas is satisfied. This ellipse is known as the apparent ellipse, and is the projection of the actual elliptical orbit of the secondary with respect to the primary on the plane of the sky. From this projected ellipse the complete elements of the orbit may be computed, where the semi-major axis can only be expressed in angular units unless the stellar parallax, and hence the distance, of the system is known.
Spectroscopic binaries
Sometimes, the only evidence of a binary star comes from the Doppler effect on its emitted light. In these cases, the binary consists of a pair of stars where the spectral lines in the light emitted from each star shifts first towards the blue, then towards the red, as each moves first towards us, and then away from us, during its motion about their common center of mass, with the period of their common orbit.
In these systems, the separation between the stars is usually very small, and the orbital velocity very high. Unless the plane of the orbit happens to be perpendicular to the line of sight, the orbital velocities have components in the line of sight, and the observed radial velocity of the system varies periodically. Since radial velocity can be measured with a spectrometer by observing the Doppler shift of the stars' spectral lines, the binaries detected in this manner are known as spectroscopic binaries. Most of these cannot be resolved as a visual binary, even with telescopes of the highest existing resolving power.
In some spectroscopic binaries, spectral lines from both stars are visible, and the lines are alternately double and single. Such a system is known as a double-lined spectroscopic binary (often denoted "SB2"). In other systems, the spectrum of only one of the stars is seen, and the lines in the spectrum shift periodically towards the blue, then towards red and back again. Such stars are known as single-lined spectroscopic binaries ("SB1").
The orbit of a spectroscopic binary is determined by making a long series of observations of the radial velocity of one or both components of the system. The observations are plotted against time, and from the resulting curve a period is determined. If the orbit is circular, then the curve is a sine curve. If the orbit is elliptical, the shape of the curve depends on the eccentricity of the ellipse and the orientation of the major axis with reference to the line of sight.
It is impossible to determine individually the semi-major axis a and the inclination of the orbit plane i. However, the product of the semi-major axis and the sine of the inclination (i.e. ) may be determined directly in linear units (e.g. kilometres). If either a or i can be determined by other means, as in the case of eclipsing binaries, a complete solution for the orbit can be found.
Binary stars that are both visual and spectroscopic binaries are rare and are a valuable source of information when found. About 40 are known. Visual binary stars often have large true separations, with periods measured in decades to centuries; consequently, they usually have orbital speeds too small to be measured spectroscopically. Conversely, spectroscopic binary stars move fast in their orbits because they are close together, usually too close to be detected as visual binaries. Binaries that are found to be both visual and spectroscopic thus must be relatively close to Earth.
Eclipsing binaries
An eclipsing binary star is a binary star system in which the orbital plane of the two stars lies so nearly in the line of sight of the observer that the components undergo mutual eclipses. In the case where the binary is also a spectroscopic binary and the parallax of the system is known, the binary is quite valuable for stellar analysis. Algol, a triple star system in the constellation Perseus, contains the best-known example of an eclipsing binary.
Eclipsing binaries are variable stars, not because the light of the individual components vary but because of the eclipses. The light curve of an eclipsing binary is characterized by periods of practically constant light, with periodic drops in intensity when one star passes in front of the other. The brightness may drop twice during the orbit, once when the secondary passes in front of the primary and once when the primary passes in front of the secondary. The deeper of the two eclipses is called the primary regardless of which star is being occulted, and if a shallow second eclipse also occurs it is called the secondary eclipse. The size of the brightness drops depends on the relative brightness of the two stars, the proportion of the occulted star that is hidden, and the surface brightness (i.e. effective temperature) of the stars. Typically the occultation of the hotter star causes the primary eclipse.
An eclipsing binary's period of orbit may be determined from a study of its light curve, and the relative sizes of the individual stars can be determined in terms of the radius of the orbit, by observing how quickly the brightness changes as the disc of the nearest star slides over the disc of the other star. If it is also a spectroscopic binary, the orbital elements can also be determined, and the mass of the stars can be determined relatively easily, which means that the relative densities of the stars can be determined in this case.
Since about 1995, measurement of extragalactic eclipsing binaries' fundamental parameters has become possible with 8-meter class telescopes. This makes it feasible to use them to directly measure the distances to external galaxies, a process that is more accurate than using standard candles. By 2006, they had been used to give direct distance estimates to the LMC, SMC, Andromeda Galaxy, and Triangulum Galaxy. Eclipsing binaries offer a direct method to gauge the distance to galaxies to an improved 5% level of accuracy.
Non-eclipsing binaries that can be detected through photometry
Nearby non-eclipsing binaries can also be photometrically detected by observing how the stars affect each other in three ways. The first is by observing extra light which the stars reflect from their companion. Second is by observing ellipsoidal light variations which are caused by deformation of the star's shape by their companions. The third method is by looking at how relativistic beaming affects the apparent magnitude of the stars. Detecting binaries with these methods requires accurate photometry.
Astrometric binaries
Astronomers have discovered some stars that seemingly orbit around an empty space. Astrometric binaries are relatively nearby stars which can be seen to wobble around a point in space, with no visible companion. The same mathematics used for ordinary binaries can be applied to infer the mass of the missing companion. The companion could be very dim, so that it is currently undetectable or masked by the glare of its primary, or it could be an object that emits little or no electromagnetic radiation, for example a neutron star.
The visible star's position is carefully measured and detected to vary, due to the gravitational influence from its counterpart. The position of the star is repeatedly measured relative to more distant stars, and then checked for periodic shifts in position. Typically this type of measurement can only be performed on nearby stars, such as those within 10 parsecs. Nearby stars often have a relatively high proper motion, so astrometric binaries will appear to follow a wobbly path across the sky.
If the companion is sufficiently massive to cause an observable shift in position of the star, then its presence can be deduced. From precise astrometric measurements of the movement of the visible star over a sufficiently long period of time, information about the mass of the companion and its orbital period can be determined. Even though the companion is not visible, the characteristics of the system can be determined from the observations using Kepler's laws.
This method of detecting binaries is also used to locate extrasolar planets orbiting a star. However, the requirements to perform this measurement are very exacting, due to the great difference in the mass ratio, and the typically long period of the planet's orbit. Detection of position shifts of a star is a very exacting science, and it is difficult to achieve the necessary precision. Space telescopes can avoid the blurring effect of Earth's atmosphere, resulting in more precise resolution.
Configuration of the system
Another classification is based on the distance between the stars, relative to their sizes:
Detached binaries are binary stars where each component is within its Roche lobe, i.e. the area where the gravitational pull of the star itself is larger than that of the other component. While on the main sequence the stars have no major effect on each other, and essentially evolve separately. Most binaries belong to this class.
Semidetached binary stars are binary stars where one of the components fills the binary star's Roche lobe and the other does not. In this interacting binary star, gas from the surface of the Roche-lobe-filling component (donor) is transferred to the other, accreting star. The mass transfer dominates the evolution of the system. In many cases, the inflowing gas forms an accretion disc around the accretor.
A contact binary is a type of binary star in which both components of the binary fill their Roche lobes. The uppermost part of the stellar atmospheres forms a common envelope that surrounds both stars. As the friction of the envelope brakes the orbital motion, the stars may eventually merge. W Ursae Majoris is an example.
Cataclysmic variables and X-ray binaries
When a binary system contains a compact object such as a white dwarf, neutron star or black hole, gas from the other (donor) star can accrete onto the compact object. This releases gravitational potential energy, causing the gas to become hotter and emit radiation. Cataclysmic variable stars, where the compact object is a white dwarf, are examples of such systems. In X-ray binaries, the compact object can be either a neutron star or a black hole. These binaries are classified as low-mass or high-mass according to the mass of the donor star. High-mass X-ray binaries contain a young, early-type, high-mass donor star which transfers mass by its stellar wind, while low-mass X-ray binaries are semidetached binaries in which gas from a late-type donor star or a white dwarf overflows the Roche lobe and falls towards the neutron star or black hole. Probably the best known example of an X-ray binary is the high-mass X-ray binary Cygnus X-1. In Cygnus X-1, the mass of the unseen companion is estimated to be about nine times that of the Sun, far exceeding the Tolman–Oppenheimer–Volkoff limit for the maximum theoretical mass of a neutron star. It is therefore believed to be a black hole; it was the first object for which this was widely believed.
Orbital period
Orbital periods can be less than an hour (for AM CVn stars), or a few days (components of Beta Lyrae), but also hundreds of thousands of years (Proxima Centauri around Alpha Centauri AB).
Variations in period
The Applegate mechanism explains long term orbital period variations seen in certain eclipsing binaries. As a main-sequence star goes through an activity cycle, the outer layers of the star are subject to a magnetic torque changing the distribution of angular momentum, resulting in a change in the star's oblateness. The orbit of the stars in the binary pair is gravitationally coupled to their shape changes, so that the period shows modulations (typically on the order of ∆P/P ~ 10−5) on the same time scale as the activity cycles (typically on the order of decades).
Another phenomenon observed in some Algol binaries has been monotonic period increases. This is quite distinct from the far more common observations of alternating period increases and decreases explained by the Applegate mechanism. Monotonic period increases have been attributed to mass transfer, usually (but not always) from the less massive to the more massive star
Designations
A and B
The components of binary stars are denoted by the suffixes A and B appended to the system's designation, A denoting the primary and B the secondary. The suffix AB may be used to denote the pair (for example, the binary star α Centauri AB consists of the stars α Centauri A and α Centauri B.) Additional letters, such as C, D, etc., may be used for systems with more than two stars. In cases where the binary star has a Bayer designation and is widely separated, it is possible that the members of the pair will be designated with superscripts; an example is Zeta Reticuli, whose components are ζ1 Reticuli and ζ2 Reticuli.
Discoverer designations
Double stars are also designated by an abbreviation giving the discoverer together with an index number. α Centauri, for example, was found to be double by Father Richaud in 1689, and so is designated RHD 1. These discoverer codes can be found in the Washington Double Star Catalog.
Hot and cold
The secondary star in a binary star system may be designated as the hot companion or cool companion, depending on its temperature relative to the primary star.
Examples:
Antares (Alpha Scorpii) is a red supergiant star in a binary system with a hotter blue main-sequence star Antares B. Antares B can therefore be termed a hot companion of the cool supergiant.
Symbiotic stars, such as R Aquarii, are binary star systems composed of a late-type giant star and a hotter companion object. Since the nature of the companion is not well-established in all cases, it may be termed a "hot companion".
The luminous blue variable Eta Carinae has been determined to be a binary star system. The secondary appears to have a higher temperature than the primary and has therefore been described as being the "hot companion" star. It may be a Wolf–Rayet star.
NASA's Kepler mission has discovered examples of eclipsing binary stars where the secondary is the hotter component. KOI-74b is a 12,000 K white dwarf companion of KOI-74 (), a 9,400 K early A-type main-sequence star. KOI-81b is a 13,000 K white dwarf companion of KOI-81 (), a 10,000 K late B-type main-sequence star.
Evolution
Formation
While it is not impossible that some binaries might be created through gravitational capture between two single stars, given the very low likelihood of such an event (three objects being actually required, as conservation of energy rules out a single gravitating body capturing another) and the high number of binaries currently in existence, this cannot be the primary formation process. The observation of binaries consisting of stars not yet on the main sequence supports the theory that binaries develop during star formation. Fragmentation of the molecular cloud during the formation of protostars is an acceptable explanation for the formation of a binary or multiple star system.
The outcome of the three-body problem, in which the three stars are of comparable mass, is that eventually one of the three stars will be ejected from the system and, assuming no significant further perturbations, the remaining two will form a stable binary system.
Mass transfer and accretion
As a main-sequence star increases in size during its evolution, it may at some point exceed its Roche lobe, meaning that some of its matter ventures into a region where the gravitational pull of its companion star is larger than its own. The result is that matter will transfer from one star to another through a process known as Roche lobe overflow (RLOF), either being absorbed by direct impact or through an accretion disc. The mathematical point through which this transfer happens is called the first Lagrangian point. It is not uncommon that the accretion disc is the brightest (and thus sometimes the only visible) element of a binary star.
If a star grows outside of its Roche lobe too fast for all abundant matter to be transferred to the other component, it is also possible that matter will leave the system through other Lagrange points or as stellar wind, thus being effectively lost to both components.
Since the evolution of a star is determined by its mass, the process influences the evolution of both companions, and creates stages that cannot be attained by single stars.
Studies of the eclipsing ternary Algol led to the Algol paradox in the theory of stellar evolution: although components of a binary star form at the same time, and massive stars evolve much faster than the less massive ones, it was observed that the more massive component Algol A is still in the main sequence, while the less massive Algol B is a subgiant at a later evolutionary stage. The paradox can be solved by mass transfer: when the more massive star became a subgiant, it filled its Roche lobe, and most of the mass was transferred to the other star, which is still in the main sequence. In some binaries similar to Algol, a gas flow can actually be seen.
Runaways and novae
It is also possible for widely separated binaries to lose gravitational contact with each other during their lifetime, as a result of external perturbations. The components will then move on to evolve as single stars. A close encounter between two binary systems can also result in the gravitational disruption of both systems, with some of the stars being ejected at high velocities, leading to runaway stars.
If a white dwarf has a close companion star that overflows its Roche lobe, the white dwarf will steadily accrete gases from the star's outer atmosphere. These are compacted on the white dwarf's surface by its intense gravity, compressed and heated to very high temperatures as additional material is drawn in. The white dwarf consists of degenerate matter and so is largely unresponsive to heat, while the accreted hydrogen is not. Hydrogen fusion can occur in a stable manner on the surface through the CNO cycle, causing the enormous amount of energy liberated by this process to blow the remaining gases away from the white dwarf's surface. The result is an extremely bright outburst of light, known as a nova.
In extreme cases this event can cause the white dwarf to exceed the Chandrasekhar limit and trigger a supernova that destroys the entire star, another possible cause for runaways. An example of such an event is the supernova SN 1572, which was observed by Tycho Brahe. The Hubble Space Telescope recently took a picture of the remnants of this event.
Astrophysics
Binaries provide the best method for astronomers to determine the mass of a distant star. The gravitational pull between them causes them to orbit around their common center of mass. From the orbital pattern of a visual binary, or the time variation of the spectrum of a spectroscopic binary, the mass of its stars can be determined, for example with the binary mass function. In this way, the relation between a star's appearance (temperature and radius) and its mass can be found, which allows for the determination of the mass of non-binaries.
Because a large proportion of stars exist in binary systems, binaries are particularly important to our understanding of the processes by which stars form. In particular, the period and masses of the binary tell us about the amount of angular momentum in the system. Because this is a conserved quantity in physics, binaries give us important clues about the conditions under which the stars were formed.
Calculating the center of mass in binary stars
In a simple binary case, the distance r1 from the center of the first star to the center of mass or barycenter is given by
where
a is the distance between the two stellar centers, and
m1 and m2 are the masses of the two stars.
If a is taken to be the semimajor axis of the orbit of one body around the other, then r1 is the semimajor axis of the first body's orbit around the center of mass or barycenter, and is the semimajor axis of the second body's orbit. When the center of mass is located within the more massive body, that body appears to wobble rather than following a discernible orbit.
Center-of-mass animations
The red cross marks the center of mass of the system. These images do not represent any specific real system.
Research findings
It is estimated that approximately one third of the star systems in the Milky Way are binary or multiple, with the remaining two thirds being single stars. The overall multiplicity frequency of ordinary stars is a monotonically increasing function of stellar mass. That is, the likelihood of being in a binary or a multi-star system steadily increases as the masses of the components increase.
There is a direct correlation between the period of revolution of a binary star and the eccentricity of its orbit, with systems of short period having smaller eccentricity. Binary stars may be found with any conceivable separation, from pairs orbiting so closely that they are practically in contact with each other, to pairs so distantly separated that their connection is indicated only by their common proper motion through space. Among gravitationally bound binary star systems, there exists a so-called log normal distribution of periods, with the majority of these systems orbiting with a period of about 100 years. This is supporting evidence for the theory that binary systems are formed during star formation.
In pairs where the two stars are of equal brightness, they are also of the same spectral type.
In systems where the brightnesses are different, the fainter star is bluer if the brighter star is a giant star, and redder if the brighter star belongs to the main sequence.
The mass of a star can be directly determined only from its gravitational attraction. Apart from the Sun and stars which act as gravitational lenses, this can be done only in binary and multiple star systems, making the binary stars an important class of stars. In the case of a visual binary star, after the orbit and the stellar parallax of the system has been determined, the combined mass of the two stars may be obtained by a direct application of the Keplerian harmonic law.
Unfortunately, it is impossible to obtain the complete orbit of a spectroscopic binary unless it is also a visual or an eclipsing binary, so from these objects only a determination of the joint product of mass and the sine of the angle of inclination relative to the line of sight is possible. In the case of eclipsing binaries which are also spectroscopic binaries, it is possible to find a complete solution for the specifications (mass, density, size, luminosity, and approximate shape) of both members of the system.
Planets
While a number of binary star systems have been found to harbor extrasolar planets, such systems are comparatively rare compared to single star systems. Observations by the Kepler space telescope have shown that most single stars of the same type as the Sun have plenty of planets, but only one-third of binary stars do. According to theoretical simulations, even widely separated binary stars often disrupt the discs of rocky grains from which protoplanets form. On the other hand, other simulations suggest that the presence of a binary companion can actually improve the rate of planet formation within stable orbital zones by "stirring up" the protoplanetary disk, increasing the accretion rate of the protoplanets within.
Detecting planets in multiple star systems introduces additional technical difficulties, which may be why they are only rarely found. Examples include the white dwarf-pulsar binary PSR B1620-26, the subgiant-red dwarf binary Gamma Cephei, and the white dwarf-red dwarf binary NN Serpentis, among others.
A study of fourteen previously known planetary systems found three of these systems to be binary systems. All planets were found to be in S-type orbits around the primary star. In these three cases the secondary star was much dimmer than the primary and so was not previously detected. This discovery resulted in a recalculation of parameters for both the planet and the primary star.
Science fiction has often featured planets of binary or ternary stars as a setting, for example, George Lucas' Tatooine from Star Wars, and one notable story, "Nightfall", even takes this to a six-star system. In reality, some orbital ranges are impossible for dynamical reasons (the planet would be expelled from its orbit relatively quickly, being either ejected from the system altogether or transferred to a more inner or outer orbital range), whilst other orbits present serious challenges for eventual biospheres because of likely extreme variations in surface temperature during different parts of the orbit. Planets that orbit just one star in a binary system are said to have "S-type" orbits, whereas those that orbit around both stars have "P-type" or "circumbinary" orbits. It is estimated that 50–60% of binary systems are capable of supporting habitable terrestrial planets within stable orbital ranges.
Examples
The large distance between the components, as well as their difference in color, make Albireo one of the easiest observable visual binaries. The brightest member, which is the third-brightest star in the constellation Cygnus, is actually a close binary itself. Also in the Cygnus constellation is Cygnus X-1, an X-ray source considered to be a black hole. It is a high-mass X-ray binary, with the optical counterpart being a variable star. Sirius is another binary and the brightest star in the night time sky, with a visual apparent magnitude of −1.46. It is located in the constellation Canis Major. In 1844 Friedrich Bessel deduced that Sirius was a binary. In 1862 Alvan Graham Clark discovered the companion (Sirius B; the visible star is Sirius A). In 1915 astronomers at the Mount Wilson Observatory determined that Sirius B was a white dwarf, the first to be discovered. In 2005, using the Hubble Space Telescope, astronomers determined Sirius B to be in diameter, with a mass that is 98% of the Sun.
An example of an eclipsing binary is Epsilon Aurigae in the constellation Auriga. The visible component belongs to the spectral class F0, the other (eclipsing) component is not visible. The last such eclipse occurred from 2009 to 2011, and it is hoped that the extensive observations that will likely be carried out may yield further insights into the nature of this system. Another eclipsing binary is Beta Lyrae, which is a semidetached binary star system in the constellation of Lyra.
Other interesting binaries include 61 Cygni (a binary in the constellation Cygnus, composed of two K class (orange) main-sequence stars, 61 Cygni A and 61 Cygni B, which is known for its large proper motion), Procyon (the brightest star in the constellation Canis Minor and the eighth-brightest star in the night time sky, which is a binary consisting of the main star with a faint white dwarf companion), SS Lacertae (an eclipsing binary which stopped eclipsing), V907 Sco (an eclipsing binary which stopped, restarted, then stopped again), BG Geminorum (an eclipsing binary which is thought to contain a black hole with a K0 star in orbit around it), and 2MASS J18082002−5104378 (a binary in the "thin disk" of the Milky Way, and containing one of the oldest known stars).
Multiple-star examples
Systems with more than two stars are termed multiple stars. Algol is the most noted ternary (long thought to be a binary), located in the constellation Perseus. Two components of the system eclipse each other, the variation in the intensity of Algol first being recorded in 1670 by Geminiano Montanari. The name Algol means "demon star" (from al-ghūl), which was probably given due to its peculiar behavior. Another visible ternary is Alpha Centauri, in the southern constellation of Centaurus, which contains the third-brightest star in the night sky, with an apparent visual magnitude of −0.01. This system also underscores the fact that no search for habitable planets is complete if binaries are discounted. Alpha Centauri A and B have an 11 AU distance at closest approach, and both should have stable habitable zones.
There are also examples of systems beyond ternaries: Castor is a sextuple star system, which is the second-brightest star in the constellation Gemini and one of the brightest stars in the nighttime sky. Astronomically, Castor was discovered to be a visual binary in 1719. Each of the components of Castor is itself a spectroscopic binary. Castor also has a faint and widely separated companion, which is also a spectroscopic binary. The Alcor–Mizar visual binary in Ursa Majoris also consists of six stars: four comprising Mizar and two comprising Alcor. QZ Carinae is a complex multiple star system made up of at least nine individual stars.
| Physical sciences | Stellar astronomy | null |
52742 | https://en.wikipedia.org/wiki/Desktop%20computer | Desktop computer | A desktop computer, often abbreviated as desktop, is a personal computer designed for regular use at a stationary location on or near a desk (as opposed to a portable computer) due to its size and power requirements. The most common configuration has a case that houses the power supply, motherboard (a printed circuit board with a microprocessor as the central processing unit, memory, bus, certain peripherals and other electronic components), disk storage (usually one or more hard disk drives, solid-state drives, optical disc drives, and in early models floppy disk drives); a keyboard and mouse for input; and a monitor, speakers, and, often, a printer for output. The case may be oriented horizontally or vertically and placed either underneath, beside, or on top of a desk.
Desktop computers with their cases oriented vertically are referred to as towers. As the majority of cases offered since the mid 1990s are in this form factor, the term desktop has been retronymically used to refer to modern cases offered in the traditional horizontal orientation.
History
Origins
Prior to the widespread use of microprocessors, a computer that could fit on a desk was considered remarkably small; the type of computers most commonly used were minicomputers, which, despite the name, were rather large and were "mini" only compared to the so-called "big iron". Early computers, and later the general purpose high throughput "mainframes", took up the space of a whole room. Minicomputers, on the contrary, generally fit into one or a few refrigerator-sized racks, or, for the few smaller ones, built into a fairly large desk, not put on top of it.
It was not until the 1970s when fully programmable computers appeared that could fit entirely on top of a desk. 1970 saw the introduction of the Datapoint 2200, a "smart" computer terminal complete with keyboard and monitor, was designed to connect with a mainframe computer but that did not stop owners from using its built-in computational abilities as a stand-alone desktop computer. The HP 9800 series, which started out as programmable calculators in 1971 but was programmable in BASIC by 1972, used a smaller version of a minicomputer design based on ROM memory and had small one-line LED alphanumeric displays and displayed graphics with a plotter. The Wang 2200 of 1973 had a full-size cathode-ray tube (CRT) and cassette tape storage. The IBM 5100 in 1975 had a small CRT display and could be programmed in BASIC and APL. These were generally expensive specialized computers sold for business or scientific uses.
Growth and development
Apple II, TRS-80 and Commodore PET were first generation personal home computers launched in 1977, which were aimed at the consumer market – rather than businessmen or computer hobbyists. Byte magazine referred to these three as the "1977 Trinity" of personal computing. Throughout the 1980s and 1990s, desktop computers became the predominant type, the most popular being the IBM PC and its clones, followed by the Apple Macintosh, with the third-placed Commodore Amiga having some success in the mid-1980s but declining by the early 1990s.
Early personal computers, like the original IBM Personal Computer, were enclosed in a "desktop case", horizontally oriented to have the display screen placed on top, thus saving space on the user's actual desk, although these cases had to be sturdy enough to support the weight of CRT displays that were widespread at the time. Over the course of the 1990s, desktop cases gradually became less common than the more-accessible tower cases that may be located on the floor under or beside a desk rather than on a desk. Not only do these tower cases have more room for expansion, they have also freed up desk space for monitors which were becoming larger every year. Desktop cases, particularly the compact form factors, remain popular for corporate computing environments and kiosks. Some computer cases can be interchangeably positioned either horizontally (desktop) or upright (mini-tower).
Influential games such as Doom and Quake during the 1990s had pushed gamers and enthusiasts to frequently upgrade to the latest CPUs and graphics cards (3dfx, ATI, and Nvidia) for their desktops (usually a tower case) in order to run these applications, though this has slowed since the late 2000s as the growing popularity of Intel integrated graphics forced game developers to scale back. Creative Technology's Sound Blaster series were a de facto standard for sound cards in desktop PCs during the 1990s until the early 2000s, when they were reduced to a niche product, as OEM desktop PCs came with sound boards integrated directly onto the motherboard.
Decline
While desktops have long been the most common configuration for PCs, by the mid-2000s the growth shifted from desktops to laptops. Notably, while desktops were mainly produced in the United States, laptops had long been produced by contract manufacturers based in Asia, such as Foxconn. This shift led to the closure of the many desktop assembly plants in the United States by 2010. Another trend around this time was the increasing proportion of inexpensive base-configuration desktops being sold, hurting PC manufacturers such as Dell whose build-to-order customization of desktops relied on upselling added features to buyers.
Battery-powered portable computers had just a 2% worldwide market share in 1986. However, laptops have become increasingly popular, both for business and personal use.
Around 109 million notebook PCs shipped worldwide in 2007, a growth of 33% compared to 2006.
In 2008, it was estimated that 145.9 million notebooks were sold and that the number would grow in 2009 to 177.7 million. The third quarter of 2008 was the first time when worldwide notebook PC shipments exceeded desktops, with 38.6 million units versus 38.5 million units.
The sales breakdown of the Apple Macintosh has seen sales of desktop Macs staying mostly constant while being surpassed by that of Mac notebooks whose sales rate has grown considerably; seven out of ten Macs sold were laptops in 2009, a ratio projected to rise to three out of four by 2010. The change in sales of form factors is due to the desktop iMac moving from affordable G3 to upscale G4 model and subsequent releases are considered premium all-in-ones. By contrast, the MSRP of the MacBook laptop lines have dropped through successive generations such that the MacBook Air and MacBook Pro constitute the lowest price of entry to a Mac, with the exception of the even more inexpensive Mac Mini (albeit without a monitor and keyboard), and the MacBooks are the top-selling form factors of the Macintosh platform today.
The decades of development mean that most people already own desktop computers that meet their needs and have no need of buying a new one merely to keep pace with advancing technology. Notably, the successive release of new versions of Windows (Windows 95, 98, XP, Vista, 7, 8, 10 and so on) had been drivers for the replacement of PCs in the 1990s, but this slowed in the 2000s due to the poor reception of Windows Vista over Windows XP. IDC analyst Jay Chou suggested that Windows 8 actually hurt sales of PCs in 2012, as businesses decided to stick with Windows 7 rather than upgrade. Some suggested that Microsoft had acknowledged "implicitly ringing the desktop PC death knell" as Windows 8 offered little upgrade in desktop PC functionality over Windows 7; instead, Windows 8's innovations were mostly on the mobile side.
The post-PC trend saw a decline in the sales of desktop and laptop PCs. The decline was attributed to increased power and applications of alternative computing devices, namely smartphones and tablet computers. Although most people exclusively use their smartphones and tablets for more basic tasks such as social media and casual gaming, these devices have in many instances replaced a second or third PC in the household that would have performed these tasks, though most families still retain a powerful PC for serious work.
Among PC form factors, desktops remain a staple in the enterprise market but lost popularity among home buyers. PC makers and electronics retailers responded by investing their engineering and marketing resources towards laptops (initially netbooks in the late 2000s, and then the higher-performance Ultrabooks from 2011 onwards), which manufacturers believed had more potential to revive the PC market than desktops.
In April 2017, StatCounter declared a "Milestone in technology history and end of an era" with the mobile Android operating system becoming more popular than Windows (the operating system that made desktops dominant over mainframe computers). Windows is still most popular on desktops (and laptops), while smartphones (and tablets) use Android or iOS.
Resurgence
Towards the middle of the 2010s, media sources began to question the existence of the post-PC trend, at least as conventionally defined, stating that the so-called post-PC devices are just other portable forms of PCs joining traditional desktop PCs which still have their own operation areas and evolve.
Although for casual use traditional desktops and laptops have seen a decline in sales, in 2018, global PC sales experienced a resurgence, driven by the business market. Desktops remain a solid fixture in the commercial and educational sectors. According to the International Data Corporation (IDC), PC sales shot up 14.8% between 2020 and 2021 and desktop market grew faster than the laptop market in the second quarter of 2021. Total PC shipments during 2021 reached 348.8 million units, up 14.8% from 2020. This represents the highest level of shipments the PC market has seen since 2012. In addition, gaming desktops have seen a global revenue increase of 54% annually. For gaming the global market of gaming desktops, laptops, and monitors is expected to grow to 61.1 million shipments by the end of 2023, up from 42.1 million, with desktops growing from 15.1 million shipments to 19 million. PC gaming as a whole accounts for 28% of the total gaming market as of 2017. This is partially due to the increasing affordability of desktop PCs.
Types
By size
Full-size
Full-sized desktops are characterized by separate display and processing components. These components are connected to each other by cables or wireless connections. They often come in a tower form factor. These computers are easy to customize and upgrade per user requirements, e.g. by expansion card.
Early extended-size (significantly larger than mainstream ATX case) tower computers sometimes were labeled as "deskside computers", but currently this naming is quite rare.
Compact
Compact desktops are reduced in physical proportions compared to full-sized desktops. They are typically small-sized, inexpensive, low-power computers designed for basic tasks such as web browsing, accessing web-based applications, document processing, and audio/video playback. Hardware specifications and processing power are usually reduced and hence make them less appropriate for running complex or resource-intensive applications. A nettop is a notable example of a compact desktop. A laptop without a screen can functionally be used as a compact desktop, sometimes called a "slabtop".
Form factor
All-in-one
An all-in-one (AIO) desktop computer integrates the system's internal components into the same case as the display, thus occupying a smaller footprint (with fewer cables) than desktops that incorporate a tower. The All-in-one systems are rarely labeled as desktop computers.
Tower
In personal computing, a tower is a form factor of desktop computer case whose height is much greater than its width, thus having the appearance of an upstanding tower block.
Pizza box form factor
In computing, a pizza box enclosure is a design for desktop computers. Pizza box cases tend to be wide and flat, resembling pizza delivery boxes and thus the name.
Cube
Cube Workstations have a cube case enclosure to house the motherboard, PCI-E expansion cards, GPU, CPU, DRAM DIMM slots, computer cooling equipment, chipsets, I/O ports, hard disk drives, and solid-state drives.
By usage
Gaming computer
Gaming computers are desktop computers with high performance CPU, GPU, and RAM optimized for playing video games at high resolution and frame rates. Gaming computer peripheries usually include mechanical keyboards for faster response time, and a gaming computer mouse which can track higher dots per inch movement.
Home theater
These desktops are connected to home entertainment systems and typically used for amusement purpose. They come with high definition display, video graphics, surround sound and TV tuner systems to complement typical PC features.
Thin client / Internet appliance
Over time some traditional desktop computers have been replaced with thin clients utilizing off-site computing solutions like the cloud. As more services and applications are served over the internet from off-site servers, local computing needs decrease, this drives desktop computers to be smaller, cheaper, and need less powerful hardware. More applications and in some cases entire virtual desktops are moved off-site and the desktop computer runs only an operating system or a shell application while the actual content is served from a server. Thin client computers may do almost all of their computing on a virtual machine in another site. Internal, hosted virtual desktops can offer users a completely consistent experience from anywhere.
Workstation
Workstations are advanced class of personal computers designed for a user and more powerful than a regular PC but less powerful than a server in regular computing. They are capable of high-resolution and three-dimensional interfaces, and typically used to perform scientific and engineering work. Like server computers, they are often connected with other workstations. The main form-factor for this class is a Tower case, but most vendors produce compact or all-in-one low-end workstations. Most tower workstations can be converted to a rack-mount version.
Desktop server
Oriented for small business class of servers; typically entry-level server machines, with similar to workstation/gaming PC computing powers and with some mainstream servers features, but with only basic graphic abilities; and some desktop servers can be converted to workstations.
Comparison with laptops
Desktops have an advantage over laptops in that the spare parts and extensions tend to be standardized, resulting in lower prices and greater availability. For example, the size and mounting of the motherboard are standardized into ATX, microATX, BTX or other form factors. Desktops have several standardized expansion slots, like conventional PCI or PCI Express, while laptops tend to have only one mini-PCI slot and one PC Card slot (or ExpressCard slot). Procedures for assembly and disassembly of desktops tend to be simple and standardized as well. This tends not to be the case for laptops, though adding or replacing some parts, like the optical drive, hard disk, or adding an extra memory module is often quite simple. This means that a desktop computer configuration, usually a tower case, can be customized and upgraded to a greater extent than laptops. This customization has kept tower cases popular among gamers and enthusiasts.
Another advantage of the desktop is that (apart from environmental concerns) power consumption is not as critical as in laptop computers because the desktop is exclusively powered from the wall socket. Desktop computers also provide more space for cooling fans and vents to dissipate heat, allowing enthusiasts to overclock with less risk. The two large microprocessor manufacturers, Intel and AMD, have developed special CPUs for mobile computers (i.e. laptops) that consume less power and lower heat, but with lower performance levels.
Laptop computers, conversely, offer portability that desktop systems (including small form factor and all-in-one desktops) cannot due to their compact size and clamshell design. The laptop's all-in-one design provides a built-in keyboard and a pointing device (such as a touchpad) for its user and can draw on power supplied by a rechargeable battery. Laptops also commonly integrate wireless technologies like Wi-Fi, Bluetooth, and 3G, giving them a broader range of options for connecting to the internet, though this trend is changing as newer desktop computers come integrated with one or more of these technologies.
A desktop computer needs a UPS to handle electrical disturbances like short interruptions, blackouts, and spikes; achieving an on-battery time of more than 20–30 minutes for a desktop PC requires a large and expensive UPS. A laptop with a sufficiently charged battery can continue to be used for hours in case of a power outage and is not affected by short power interruptions and blackouts.
A desktop computer often has the advantage over a comparable laptop in computational capacity. Overclocking is often more feasible on a desktop than on a laptop; similarly, hardware add-ons such as discrete graphics co-processors may be possible to install only in a desktop.
| Technology | Computer hardware | null |
52778 | https://en.wikipedia.org/wiki/Sodium%20cyanide | Sodium cyanide | Sodium cyanide is a compound with the formula NaCN and the structure . It is a white, water-soluble solid. Cyanide has a high affinity for metals, which leads to the high toxicity of this salt. Its main application, in gold mining, also exploits its high reactivity toward metals. It is a moderately strong base.
Production and chemical properties
Sodium cyanide is produced by treating hydrogen cyanide with sodium hydroxide:
Worldwide production was estimated at 500,000 tons in the year 2006. Formerly it was prepared by the Castner process involving the reaction of sodium amide with carbon at elevated temperatures.
The structure of solid NaCN is related to that of sodium chloride. The anions and cations are each six-coordinate. Potassium cyanide (KCN) adopts a similar structure.
When treated with acid, it forms the toxic gas hydrogen cyanide:
Because the salt is derived from a weak acid, sodium cyanide readily reverts to HCN by hydrolysis; the moist solid emits small amounts of hydrogen cyanide, which is thought to smell like bitter almonds (not everyone can smell it—the ability thereof is due to a genetic trait). Sodium cyanide reacts rapidly with strong acids to release hydrogen cyanide. This dangerous process represents a significant risk associated with cyanide salts. It is detoxified most efficiently with hydrogen peroxide () to produce sodium cyanate (NaOCN) and water:
Applications
Cyanide mining
Gold cyanidation (also known as the cyanide process) is the dominant technique for extracting gold, much of which is obtained from low-grade ore. More than 70% of cyanide consumption globally is used for this purpose. The application exploits the high affinity of gold(I) for cyanide, which induces gold metal to oxidize and dissolve in the presence of air (oxygen) and water, producing the salt sodium dicyanoaurate (or sodium gold cyanide) ():
A similar process uses potassium cyanide (KCN, a close relative of sodium cyanide) to produce potassium dicyanoaurate ().
Chemical feedstock
Several commercially significant chemical compounds are derived from cyanide, including cyanuric chloride, cyanogen chloride, and many nitriles. In organic synthesis, cyanide, which is classified as a strong nucleophile, is used to prepare nitriles, which occur widely in many chemicals, including pharmaceuticals. Illustrative is the synthesis of benzyl cyanide by the reaction of benzyl chloride and sodium cyanide.
Niche uses
Being highly toxic, sodium cyanide is used to kill or stun rapidly such as in collecting jars used by entomologists and in widely illegal cyanide fishing.
It was used as an insecticide, rodenticide and antibacterial, but these uses were cancelled by the EPA in 1987.
Toxicity
Sodium cyanide, like other soluble cyanide salts, is among the most rapidly acting of all known poisons. NaCN is a potent inhibitor of respiration, acting on mitochondrial cytochrome oxidase and hence blocking electron transport. This results in decreased oxidative metabolism and oxygen utilization. Lactic acidosis then occurs as a consequence of anaerobic metabolism. An oral dosage as small as 200–300 mg can be fatal.
| Physical sciences | Cyanide salts | Chemistry |
52812 | https://en.wikipedia.org/wiki/Humidity | Humidity | Humidity is the concentration of water vapor present in the air. Water vapor, the gaseous state of water, is generally invisible to the human eye. Humidity indicates the likelihood for precipitation, dew, or fog to be present.
Humidity depends on the temperature and pressure of the system of interest. The same amount of water vapor results in higher relative humidity in cool air than warm air. A related parameter is the dew point. The amount of water vapor needed to achieve saturation increases as the temperature increases. As the temperature of a parcel of air decreases it will eventually reach the saturation point without adding or losing water mass. The amount of water vapor contained within a parcel of air can vary significantly. For example, a parcel of air near saturation may contain 8 g of water per cubic metre of air at , and 28 g of water per cubic metre of air at
Three primary measurements of humidity are widely employed: absolute, relative, and specific. Absolute humidity is expressed as either mass of water vapor per volume of moist air (in grams per cubic meter) or as mass of water vapor per mass of dry air (usually in grams per kilogram). Relative humidity, often expressed as a percentage, indicates a present state of absolute humidity relative to a maximum humidity given the same temperature. Specific humidity is the ratio of water vapor mass to total moist air parcel mass.
Humidity plays an important role for surface life. For animal life dependent on perspiration (sweating) to regulate internal body temperature, high humidity impairs heat exchange efficiency by reducing the rate of moisture evaporation from skin surfaces. This effect can be calculated using a heat index table, or alternatively using a similar humidex.
The notion of air "holding" water vapor or being "saturated" by it is often mentioned in connection with the concept of relative humidity. This, however, is misleading—the amount of water vapor that enters (or can enter) a given space at a given temperature is almost independent of the amount of air (nitrogen, oxygen, etc.) that is present. Indeed, a vacuum has approximately the same equilibrium capacity to hold water vapor as the same volume filled with air; both are given by the equilibrium vapor pressure of water at the given temperature. There is a very small difference described under "Enhancement factor" below, which can be neglected in many calculations unless great accuracy is required.
Definitions
Absolute humidity
Absolute humidity is the total mass of water vapor (gas form of water) present in a given volume or mass of air. It does not take temperature into consideration. Absolute humidity in the atmosphere ranges from near zero to roughly per cubic metre when the air is saturated at .
Air is a gas, and its volume varies with pressure and temperature, per Boyles law. Absolute humidity is defined as water mass per volume of air. A given mass of air will grow or shrink as the temperature or pressure varies. So the absolute humidity of a mass of air will vary due to changes in temperature or pressure, even when the proportion of water in that mass of air (its specific humidity) remains constant. This makes the term absolute humidity as defined not ideal for some situations.
Absolute humidity is the mass of the water vapor , divided by the volume of the air and water vapor mixture , which can be expressed as:
In the equation above, if the volume is not set, the absolute humidity varies with changes in air temperature or pressure. Because of this variability, use of the term absolute humidity as defined is inappropriate for computations in chemical engineering, such as drying, where temperature variations might be significant. As a result, absolute humidity in chemical engineering may refer to mass of water vapor per unit mass of dry air, also known as the humidity ratio or mass mixing ratio (see "specific humidity" below), which is better suited for heat and mass balance calculations. Mass of water per unit volume as in the equation above is also defined as volumetric humidity. Because of the potential confusion, British Standard BS 1339 suggests avoiding the term "absolute humidity". Units should always be carefully checked. Many humidity charts are given in g/kg or kg/kg, but any mass units may be used.
Relative humidity
Relative humidity is the ratio of how much water vapour is in the air to how much water vapour the air could potentially contain at a given temperature. It varies with the temperature of the air: colder air can contain less vapour, and water will tend to condense out of the air more at lower temperatures. So changing the temperature of air can change the relative humidity, even when the absolute humidity remains constant.
Chilling air increases the relative humidity. If the relative humidity rises over 100% (the dew point) and there is an available surface or particle, the water vapour will condense into liquid or ice. Likewise, warming air decreases the relative humidity. Warming some air containing a fog may cause that fog to evaporate, as the droplets are prone to total evaporation due to the lowering partial pressure of water vapour in that air, as the temperature rises.
Relative humidity only considers the invisible water vapour. Mists, clouds, fogs and aerosols of water do not count towards the measure of relative humidity of the air, although their presence is an indication that a body of air may be close to the dew point.
Relative humidity is normally expressed as a percentage; a higher percentage means that the air–water mixture is more humid. At 100% relative humidity, the air is saturated and is at its dew point. In the absence of a foreign body on which droplets or crystals can nucleate, the relative humidity can exceed 100%, in which case the air is said to be supersaturated. Introduction of some particles or a surface to a body of air above 100% relative humidity will allow condensation or ice to form on those nuclei, thereby removing some of the vapour and lowering the humidity.
In a scientific notion, the relative humidity ( or ) of an air-water mixture is defined as the ratio of the partial pressure of water vapor () in air to the saturation vapor pressure () of water at the same temperature, usually expressed as a percentage:
Relative humidity is an important metric used in weather forecasts and reports, as it is an indicator of the likelihood of precipitation, dew, or fog. In hot summer weather, a rise in relative humidity increases the apparent temperature to humans (and other animals) by hindering the evaporation of perspiration from the skin. For example, according to the heat index, a relative humidity of 75% at air temperature of would feel like .
Relative humidity is also a key metric used to evaluate when it is appropriate to install flooring over a concrete slab.
Specific humidity
Specific humidity (or moisture content) is the ratio of the mass of water vapor to the total mass of the air parcel. Specific humidity is approximately equal to the mixing ratio, which is defined as the ratio of the mass of water vapor in an air parcel to the mass of dry air for the same parcel. It is typically represented with the symbol ω, and is commonly used in HVAC system design.
Related concepts
The term relative humidity is reserved for systems of water vapor in air. The term relative saturation is used to describe the analogous property for systems consisting of a condensable phase other than water in a non-condensable phase other than air.
Measurement
A device used to measure humidity of air is called a psychrometer or hygrometer. A humidistat is a humidity-triggered switch, often used to control a humidifier or a dehumidifier.
The humidity of an air and water vapor mixture is determined through the use of psychrometric charts if both the dry bulb temperature (T) and the wet bulb temperature (Tw) of the mixture are known. These quantities are readily estimated by using a sling psychrometer.
There are several empirical formulas that can be used to estimate the equilibrium vapor pressure of water vapor as a function of temperature. The Antoine equation is among the least complex of these, having only three parameters (A, B, and C). Other formulas, such as the Goff–Gratch equation and the Magnus–Tetens approximation, are more complicated but yield better accuracy.
The Arden Buck equation is commonly encountered in the literature regarding this topic:
where is the dry-bulb temperature expressed in degrees Celsius (°C), is the absolute pressure expressed in millibars, and is the equilibrium vapor pressure expressed in millibars. Buck has reported that the maximal relative error is less than 0.20% between when this particular form of the generalized formula is used to estimate the equilibrium vapor pressure of water.
There are various devices used to measure and regulate humidity. Calibration standards for the most accurate measurement include the gravimetric hygrometer, chilled mirror hygrometer, and electrolytic hygrometer. The gravimetric method, while the most accurate, is very cumbersome. For fast and very accurate measurement the chilled mirror method is effective. For process on-line measurements, the most commonly used sensors nowadays are based on capacitance measurements to measure relative humidity, frequently with internal conversions to display absolute humidity as well. These are cheap, simple, generally accurate and relatively robust. All humidity sensors face problems in measuring dust-laden gas, such as exhaust streams from clothes dryers.
Humidity is also measured on a global scale using remotely placed satellites. These satellites are able to detect the concentration of water in the troposphere at altitudes between . Satellites that can measure water vapor have sensors that are sensitive to infrared radiation. Water vapor specifically absorbs and re-radiates radiation in this spectral band. Satellite water vapor imagery plays an important role in monitoring climate conditions (like the formation of thunderstorms) and in the development of weather forecasts.
Air density and volume
Humidity depends on water vaporization and condensation, which, in turn, mainly depends on temperature. Therefore, when applying more pressure to a gas saturated with water, all components will initially decrease in volume approximately according to the ideal gas law. However, some of the water will condense until returning to almost the same humidity as before, giving the resulting total volume deviating from what the ideal gas law predicted.
Conversely, decreasing temperature would also make some water condense, again making the final volume deviate from predicted by the ideal gas law. Therefore, gas volume may alternatively be expressed as the dry volume, excluding the humidity content. This fraction more accurately follows the ideal gas law. On the contrary the saturated volume is the volume a gas mixture would have if humidity was added to it until saturation (or 100% relative humidity).
Humid air is less dense than dry air because a molecule of water () is less massive than either a molecule of nitrogen () or a molecule of oxygen (). About 78% of the molecules in dry air are nitrogen (N2). Another 21% of the molecules in dry air are oxygen (O2). The final 1% of dry air is a mixture of other gases.
For any gas, at a given temperature and pressure, the number of molecules present in a particular volume is constant. Therefore, when some number N of water molecules (vapor) is introduced into a volume of dry air, the number of air molecules in that volume must decrease by the same number N for the pressure to remain constant without using a change in temperature. The numbers are exactly equal if we consider the gases as ideal. The addition of water molecules, or any other molecules, to a gas, without removal of an equal number of other molecules, will necessarily require a change in temperature, pressure, or total volume; that is, a change in at least one of these three parameters.
If temperature and pressure remain constant, the volume increases, and the dry air molecules that were displaced will initially move out into the additional volume, after which the mixture will eventually become uniform through diffusion. Hence the mass per unit volume of the gas—its density—decreases. Isaac Newton discovered this phenomenon and wrote about it in his book Opticks.
Pressure dependence
The relative humidity of an air–water system is dependent not only on the temperature but also on the absolute pressure of the system of interest. This dependence is demonstrated by considering the air–water system shown below. The system is closed (i.e., no matter enters or leaves the system).
If the system at State A is isobarically heated (heating with no change in system pressure), then the relative humidity of the system decreases because the equilibrium vapor pressure of water increases with increasing temperature. This is shown in State B.
If the system at State A is isothermally compressed (compressed with no change in system temperature), then the relative humidity of the system increases because the partial pressure of water in the system increases with the volume reduction. This is shown in State C. Above 202.64 kPa, the RH would exceed 100% and water may begin to condense.
If the pressure of State A was changed by simply adding more dry air, without changing the volume, the relative humidity would not change.
Therefore, a change in relative humidity can be explained by a change in system temperature, a change in the volume of the system, or change in both of these system properties.
Enhancement factor
The enhancement factor is defined as the ratio of the saturated vapor pressure of water in moist air to the saturated vapor pressure of pure water:
The enhancement factor is equal to unity for ideal gas systems. However, in real systems the interaction effects between gas molecules result in a small increase of the equilibrium vapor pressure of water in air relative to equilibrium vapor pressure of pure water vapor. Therefore, the enhancement factor is normally slightly greater than unity for real systems.
The enhancement factor is commonly used to correct the equilibrium vapor pressure of water vapor when empirical relationships, such as those developed by Wexler, Goff, and Gratch, are used to estimate the properties of psychrometric systems.
Buck has reported that, at sea level, the vapor pressure of water in saturated moist air amounts to an increase of approximately 0.5% over the equilibrium vapor pressure of pure water.
Effects
Climate control refers to the control of temperature and relative humidity in buildings, vehicles and other enclosed spaces for the purpose of providing for human comfort, health and safety, and of meeting environmental requirements of machines, sensitive materials (for example, historic) and technical processes.
Climate
While humidity itself is a climate variable, it also affects other climate variables. Environmental humidity is affected by winds and by rainfall.
The most humid cities on Earth are generally located closer to the equator, near coastal regions. Cities in parts of Asia and Oceania are among the most humid. Bangkok, Ho Chi Minh City, Kuala Lumpur, Hong Kong, Manila, Jakarta, Naha, Singapore, Kaohsiung and Taipei have very high humidity most or all year round because of their proximity to water bodies and the equator and often overcast weather.
Some places experience extreme humidity during their rainy seasons combined with warmth giving the feel of a lukewarm sauna, such as Kolkata, Chennai and Kochi in India, and Lahore in Pakistan. Sukkur city located on the Indus River in Pakistan has some of the highest and most uncomfortable dew points in the country, frequently exceeding in the monsoon season.
High temperatures combine with the high dew point to create heat index in excess of . Darwin experiences an extremely humid wet season from December to April. Houston, Miami, San Diego, Osaka, Shanghai, Shenzhen and Tokyo also have an extreme humid period in their summer months. During the South-west and North-east Monsoon seasons (respectively, late May to September and November to March), expect heavy rains and a relatively high humidity post-rainfall.
Outside the monsoon seasons, humidity is high (in comparison to countries further from the Equator), but completely sunny days abound. In cooler places such as Northern Tasmania, Australia, high humidity is experienced all year due to the ocean between mainland Australia and Tasmania. In the summer the hot dry air is absorbed by this ocean and the temperature rarely climbs above .
Global climate
Humidity affects the energy budget and thereby influences temperatures in two major ways. First, water vapor in the atmosphere contains "latent" energy. During transpiration or evaporation, this latent heat is removed from surface liquid, cooling the Earth's surface. This is the biggest non-radiative cooling effect at the surface. It compensates for roughly 70% of the average net radiative warming at the surface.
Second, water vapor is the most abundant of all greenhouse gases. Water vapor, like a green lens that allows green light to pass through it but absorbs red light, is a "selective absorber". Like the other greenhouse gasses, water vapor is transparent to most solar energy. However, it absorbs the infrared energy emitted (radiated) upward by the Earth's surface, which is the reason that humid areas experience very little nocturnal cooling but dry desert regions cool considerably at night. This selective absorption causes the greenhouse effect. It raises the surface temperature substantially above its theoretical radiative equilibrium temperature with the sun, and water vapor is the cause of more of this warming than any other greenhouse gas.
Unlike most other greenhouse gases, however, water is not merely below its boiling point in all regions of the Earth, but below its freezing point at many altitudes. As a condensible greenhouse gas, it precipitates, with a much lower scale height and shorter atmospheric lifetime — weeks instead of decades. Without other greenhouse gases, Earth's blackbody temperature, below the freezing point of water, would cause water vapor to be removed from the atmosphere. Water vapor is thus a "slave" to the non-condensible greenhouse gases.
Animal and plant life
Humidity is one of the fundamental abiotic factors that defines any habitat (the tundra, wetlands, and the desert are a few examples), and is a determinant of which animals and plants can thrive in a given environment.
The human body dissipates heat through perspiration and its evaporation. Heat convection, to the surrounding air, and thermal radiation are the primary modes of heat transport from the body. Under conditions of high humidity, the rate of evaporation of sweat from the skin decreases. Also, if the atmosphere is as warm or warmer than the skin during times of high humidity, blood brought to the body surface cannot dissipate heat by conduction to the air. With so much blood going to the external surface of the body, less goes to the active muscles, the brain, and other internal organs. Physical strength declines, and fatigue occurs sooner than it would otherwise. Alertness and mental capacity also may be affected, resulting in heat stroke or hyperthermia.
Domesticated plants and animals (e.g. lizards) require regular upkeep of humidity percent when grown in-home and container conditions, for optimal thriving environment.
Human comfort
Although humidity is an important factor for thermal comfort, humans are more sensitive to variations in temperature than they are to changes in relative humidity. Humidity has a small effect on thermal comfort outdoors when air temperatures are low, a slightly more pronounced effect at moderate air temperatures, and a much stronger influence at higher air temperatures.
Humans are sensitive to humid air because the human body uses evaporative cooling as the primary mechanism to regulate temperature. Under humid conditions, the rate at which perspiration evaporates on the skin is lower than it would be under arid conditions. Because humans perceive the rate of heat transfer from the body rather than temperature itself, we feel warmer when the relative humidity is high than when it is low.
Humans can be comfortable within a wide range of humidities depending on the temperature—from 30 to 70%—but ideally not above the Absolute (60 °F Dew Point), between 40% and 60%. In general, higher temperatures will require lower humidities to achieve thermal comfort compared to lower temperatures, with all other factors held constant. For example, with clothing level = 1, metabolic rate = 1.1, and air speed 0.1 m/s, a change in air temperature and mean radiant temperature from 20 °C to 24 °C would lower the maximum acceptable relative humidity from 100% to 65% to maintain thermal comfort conditions. The CBE Thermal Comfort Tool can be used to demonstrate the effect of relative humidity for specific thermal comfort conditions and it can be used to demonstrate compliance with ASHRAE Standard 55–2017.
Some people experience difficulty breathing in humid environments. Some cases may possibly be related to respiratory conditions such as asthma, while others may be the product of anxiety. Affected people will often hyperventilate in response, causing sensations of numbness, faintness, and loss of concentration, among others.
Very low humidity can create discomfort, respiratory problems, and aggravate allergies in some individuals. Low humidity causes tissue lining nasal passages to dry, crack and become more susceptible to penetration of rhinovirus cold viruses. Extremely low (below 20%) relative humidities may also cause eye irritation. The use of a humidifier in homes, especially bedrooms, can help with these symptoms. Indoor relative humidities kept above 30% reduce the likelihood of the occupant's nasal passages drying out, especially in winter.
Air conditioning reduces discomfort by reducing not just temperature but humidity as well. Heating cold outdoor air can decrease relative humidity levels indoors to below 30%. According to ASHRAE Standard 55-2017: Thermal Environmental Conditions for Human Occupancy, indoor thermal comfort can be achieved through the PMV method with relative humidities ranging from 0% to 100%, depending on the levels of the other factors contributing to thermal comfort. However, the recommended range of indoor relative humidity in air conditioned buildings is generally 30–60%.
Human health
Higher humidity reduces the infectivity of aerosolized influenza virus. A study concluded, "Maintaining indoor relative humidity >40% will significantly reduce the infectivity of aerosolized virus."
Excess moisture in buildings expose occupants to fungal spores, cell fragments, or mycotoxins. Infants in homes with mold have a much greater risk of developing asthma and allergic rhinitis. More than half of adult workers in moldy/humid buildings develop nasal or sinus symptoms due to mold exposure.
Mucociliary clearance in the respiratory tract is also hindered by low humidity. One study in dogs found that mucus transport was lower at an absolute humidity of 9 g/m3 than at 30 g/m3.
Increased humidity can also lead to changes in total body water that usually leads to moderate weight gain, especially if one is acclimated to working or exercising in hot and humid weather.
Building construction
Common construction methods often produce building enclosures with a poor thermal boundary, requiring an insulation and air barrier system designed to retain indoor environmental conditions while resisting external environmental conditions. The energy-efficient, heavily sealed architecture introduced in the 20th century also sealed off the movement of moisture, and this has resulted in a secondary problem of condensation forming in and around walls, which encourages the development of mold and mildew. Additionally, buildings with foundations not properly sealed will allow water to flow through the walls due to capillary action of pores found in masonry products. Solutions for energy-efficient buildings that avoid condensation are a current topic of architecture.
For climate control in buildings using HVAC systems, the key is to maintain the relative humidity at a comfortable range—low enough to be comfortable but high enough to avoid problems associated with very dry air.
When the temperature is high and the relative humidity is low, evaporation of water is rapid; soil dries, wet clothes hung on a line or rack dry quickly, and perspiration readily evaporates from the skin. Wooden furniture can shrink, causing the paint that covers these surfaces to fracture.
When the temperature is low and the relative humidity is high, evaporation of water is slow. When relative humidity approaches 100%, condensation can occur on surfaces, leading to problems with mold, corrosion, decay, and other moisture-related deterioration. Condensation can pose a safety risk as it can promote the growth of mold and wood rot as well as possibly freezing emergency exits shut.
Certain production and technical processes and treatments in factories, laboratories, hospitals, and other facilities require specific relative humidity levels to be maintained using humidifiers, dehumidifiers and associated control systems.
Vehicles
The basic principles for buildings, above, also apply to vehicles. In addition, there may be safety considerations. For instance, high humidity inside a vehicle can lead to problems of condensation, such as misting of windshields and shorting of electrical components. In vehicles and pressure vessels such as pressurized airliners, submersibles and spacecraft, these considerations may be critical to safety, and complex environmental control systems including equipment to maintain pressure are needed.
Aviation
Airliners operate with low internal relative humidity, often under 20%, especially on long flights. The low humidity is a consequence of drawing in the very cold air with a low absolute humidity, which is found at airliner cruising altitudes. Subsequent warming of this air lowers its relative humidity. This causes discomfort such as sore eyes, dry skin, and drying out of mucosa, but humidifiers are not employed to raise it to comfortable mid-range levels because the volume of water required to be carried on board can be a significant weight penalty. As airliners descend from colder altitudes into warmer air, perhaps even flying through clouds a few thousand feet above the ground, the ambient relative humidity can increase dramatically.
Some of this moist air is usually drawn into the pressurized aircraft cabin and into other non-pressurized areas of the aircraft and condenses on the cold aircraft skin. Liquid water can usually be seen running along the aircraft skin, both on the inside and outside of the cabin. Because of the drastic changes in relative humidity inside the vehicle, components must be qualified to operate in those environments. The recommended environmental qualifications for most commercial aircraft components is listed in RTCA DO-160.
Cold, humid air can promote the formation of ice, which is a danger to aircraft as it affects the wing profile and increases weight. Naturally aspirated internal combustion engines have a further danger of ice forming inside the carburetor. Aviation weather reports (METARs) therefore include an indication of relative humidity, usually in the form of the dew point.
Pilots must take humidity into account when calculating takeoff distances, because high humidity requires longer runways and will decrease climb performance.
Density altitude is the altitude relative to the standard atmosphere conditions (International Standard Atmosphere) at which the air density would be equal to the indicated air density at the place of observation, or, in other words, the height when measured in terms of the density of the air rather than the distance from the ground. "Density Altitude" is the pressure altitude adjusted for non-standard temperature.
An increase in temperature, and, to a much lesser degree, humidity, will cause an increase in density altitude. Thus, in hot and humid conditions, the density altitude at a particular location may be significantly higher than the true altitude.
Electronics
Electronic devices are often rated to operate only under certain humidity conditions (e.g., 10% to 90%). The optimal humidity for electronic devices is 30% to 65%. At the top end of the range, moisture may increase the conductivity of permeable insulators leading to malfunction. Too low humidity may make materials brittle. A particular danger to electronic items, regardless of the stated operating humidity range, is condensation. When an electronic item is moved from a cold place (e.g., garage, car, shed, air conditioned space in the tropics) to a warm humid place (house, outside tropics), condensation may coat circuit boards and other insulators, leading to short circuit inside the equipment. Such short circuits may cause substantial permanent damage if the equipment is powered on before the condensation has evaporated. A similar condensation effect can often be observed when a person wearing glasses comes in from the cold (i.e. the glasses become foggy).
It is advisable to allow electronic equipment to acclimatise for several hours, after being brought in from the cold, before powering on. Some electronic devices can detect such a change and indicate, when plugged in and usually with a small droplet symbol, that they cannot be used until the risk from condensation has passed. In situations where time is critical, increasing air flow through the device's internals, such as removing the side panel from a PC case and directing a fan to blow into the case, will reduce significantly the time needed to acclimatise to the new environment.
In contrast, a very low humidity level favors the build-up of static electricity, which may result in spontaneous shutdown of computers when discharges occur. Apart from spurious erratic function, electrostatic discharges can cause dielectric breakdown in solid-state devices, resulting in irreversible damage. Data centers often monitor relative humidity levels for these reasons.
Industry
High humidity can often have a negative effect on the capacity of chemical plants and refineries that use furnaces as part of a certain processes (e.g., steam reforming, wet sulfuric acid processes). For example, because humidity reduces ambient oxygen concentrations (dry air is typically 20.9% oxygen, but at 100% relative humidity the air is 20.4% oxygen), flue gas fans must intake air at a higher rate than would otherwise be required to maintain the same firing rate.
Baking
High humidity in the oven, represented by an elevated wet-bulb temperature, increases the thermal conductivity of the air around the baked item, leading to a quicker baking process or even burning. Conversely, low humidity slows the baking process down.
Other important facts
At 100% relative humidity, air is saturated and at its dew point: the water vapor pressure would permit neither evaporation of nearby liquid water nor condensation to grow the nearby water; neither sublimation of nearby ice nor deposition to grow the nearby ice.
Relative humidity can exceed 100%, in which case the air is supersaturated. Cloud formation requires supersaturated air. Cloud condensation nuclei lower the level of supersaturation required to form fogs and clouds – in the absence of nuclei around which droplets or ice can form, a higher level of supersaturation is required for these droplets or ice crystals to form spontaneously. In the Wilson cloud chamber, which is used in nuclear physics experiments, a state of supersaturation is created within the chamber, and moving subatomic particles act as condensation nuclei so trails of fog show the paths of those particles.
For a given dew point and its corresponding absolute humidity, the relative humidity will change inversely, albeit nonlinearly, with the temperature. This is because the vapor pressure of water increases with temperature—the operative principle behind everything from hair dryers to dehumidifiers.
Due to the increasing potential for a higher water vapor partial pressure at higher air temperatures, the water content of air at sea level can get as high as 3% by mass at compared to no more than about 0.5% by mass at . This explains the low levels (in the absence of measures to add moisture) of humidity in heated structures during winter, resulting in dry skin, itchy eyes, and persistence of static electric charges. Even with saturation (100% relative humidity) outdoors, heating of infiltrated outside air that comes indoors raises its moisture capacity, which lowers relative humidity and increases evaporation rates from moist surfaces indoors, including human bodies and household plants.
Similarly, during summer in humid climates a great deal of liquid water condenses from air cooled in air conditioners. Warmer air is cooled below its dew point, and the excess water vapor condenses. This phenomenon is the same as that which causes water droplets to form on the outside of a cup containing an ice-cold drink.
A useful rule of thumb is that the maximum absolute humidity doubles for every increase in temperature. Thus, the relative humidity will drop by a factor of 2 for each increase in temperature, assuming conservation of absolute moisture. For example, in the range of normal temperatures, air at and 50% relative humidity will become saturated if cooled to , its dew point, and air at 80% relative humidity warmed to will have a relative humidity of only 29% and feel dry. By comparison, thermal comfort standard ASHRAE 55 requires systems designed to control humidity to maintain a dew point of though no lower humidity limit is established.
Water vapor is a lighter gas than other gaseous components of air at the same temperature, so humid air will tend to rise by natural convection. This is a mechanism behind thunderstorms and other weather phenomena. Relative humidity is often mentioned in weather forecasts and reports, as it is an indicator of the likelihood of dew, or fog. In hot summer weather, it also increases the apparent temperature to humans (and other animals) by hindering the evaporation of perspiration from the skin as the relative humidity rises. This effect is calculated as the heat index or humidex.
A device used to measure humidity is called a hygrometer; one used to regulate it is called a humidistat, or sometimes hygrostat. These are analogous to a thermometer and thermostat for temperature, respectively.
The field concerned with the study of physical and thermodynamic properties of gas–vapor mixtures is named psychrometrics.
Relationship between absolute humidity, relative humidity, and temperature
| Physical sciences | Meteorology: General | null |
52838 | https://en.wikipedia.org/wiki/Charon%20%28moon%29 | Charon (moon) | Charon ( or ), or (134340) Pluto I, is the largest of the five known natural satellites of the dwarf planet Pluto. It has a mean radius of . Charon is the sixth-largest known trans-Neptunian object after Pluto, Eris, Haumea, Makemake, and Gonggong. It was discovered in 1978 at the United States Naval Observatory in Washington, D.C., using photographic plates taken at the United States Naval Observatory Flagstaff Station (NOFS).
With half the diameter and one-eighth the mass of Pluto, Charon is a very large moon in comparison to its parent body. Its gravitational influence is such that the barycenter of the Plutonian system lies outside Pluto, and the two bodies are tidally locked to each other. The dwarf planet systems Pluto–Charon and Eris–Dysnomia are the only known examples of mutual tidal locking in the Solar System, though it is likely that –Vanth is another.
The reddish-brown cap of the north pole of Charon is composed of tholins, organic macromolecules that may be essential ingredients of life. These tholins were produced from methane, nitrogen, and related gases which may have been released by cryovolcanic eruptions on the moon, or may have been transferred over from the atmosphere of Pluto to the orbiting moon.
The New Horizons spacecraft is the only probe that has visited the Pluto system. It approached Charon to within in 2015.
Discovery
Charon was discovered by United States Naval Observatory astronomer James Christy, using the telescope at United States Naval Observatory Flagstaff Station (NOFS). On June 22, 1978, he had been examining highly magnified images of Pluto on photographic plates taken with the telescope two months prior. Christy noticed that a slight elongation appeared periodically. The bulge was confirmed on plates dating back to April 29, 1965. Subsequent observations of Pluto determined that the bulge was due to a smaller accompanying body. The periodicity of the bulge corresponded to Pluto's rotation period, which was previously known from Pluto's light curve. This indicated a synchronous orbit, which strongly suggested that the bulge effect was real and not spurious. This resulted in reassessments of Pluto's size, mass, and other physical characteristics because the calculated mass and albedo of the Pluto–Charon system had previously been attributed to Pluto alone. The International Astronomical Union formally announced Christy's discovery to the world on July 7, 1978.
Doubts about Charon's existence were erased when it and Pluto entered a five-year period of mutual eclipses and transits between 1985 and 1990. This occurs when the Pluto–Charon orbital plane is edge-on as seen from Earth, which only happens at two intervals in Pluto's 248-year orbital period. It was fortuitous that one of these intervals happened to occur soon after Charon's discovery.
Name
Charon was first given the temporary designation S/1978 P 1, after its discovery, following the then recently instituted convention. On June 24, 1978, Christy first suggested the name Charon as a scientific-sounding version of his wife Charlene's nickname, "Char". Although colleagues at the Naval Observatory proposed Persephone, Christy stuck with Charon after discovering that it was serendipitously the name of an appropriate mythological figure: Charon (; ) is the ferryman of the dead, closely associated with the god Pluto. The IAU officially adopted the name in late 1985, and it was announced on January 3, 1986.
Coincidentally, nearly four decades before Charon's discovery, science fiction author Edmond Hamilton had invented three moons of Pluto for his 1940 novel Calling Captain Future and named them Charon, Styx, and Cerberus;
Styx and Kerberos are the two smallest Plutonian moons, and were named in 2013.
There is minor debate over the preferred pronunciation of the name. The mythological figure is pronounced with a sound, and this is often followed for the moon as well. However, Christy himself pronounced the initial as a sound, as he had named the moon after his wife Charlene. Many English-speaking astronomers follow the classical convention, but others follow Christy's, and that is the prescribed pronunciation at NASA and of the New Horizons team.
Orbit
Charon and Pluto orbit each other every 6.387 days. The two objects are gravitationally locked to one another, so each keeps the same face towards the other. This is a case of mutual tidal locking, as compared to that of the Earth and the Moon, where the Moon always shows the same face to Earth, but not vice versa. The average distance between Charon and Pluto is . The discovery of Charon allowed astronomers to calculate accurately the mass of the Plutonian system, and mutual occultations revealed their sizes. However, neither indicated the two bodies' individual masses. Those could only be estimated, until the discovery of Pluto's outer moons in late 2005. Details in the orbits of the outer moons then revealed that Charon has approximately 12% of the mass of Pluto.
Formation
Simulation work published in 2005 by Robin Canup suggested that Charon could have been formed by a collision around 4.5 billion years ago, much like Earth and the Moon. In this model, a large Kuiper belt object struck Pluto at high velocity, destroying itself and blasting off much of Pluto's outer mantle, and Charon coalesced from the debris. However, such an impact should result in an icier Charon and rockier Pluto than scientists have found. It is now thought that Pluto and Charon might have been two bodies that collided before going into orbit around each other. The collision would have been violent enough to boil off volatile ices like methane () but not violent enough to have destroyed either body. The very similar density of Pluto and Charon implies that the parent bodies were not fully differentiated when the impact occurred. The two bodies would have been stuck for a while, before separating from each other again, while remaining gravitationally bound. The internal heat in both bodies, created from both the collision and then the tidal friction as they separated, would have been sufficient to create Pluto's subsurface ocean without the need for radioactive elements.
Physical characteristics
Charon's diameter is , just over half that of Pluto. Larger than the dwarf planet Ceres, it is the twelfth-largest natural satellite in the Solar System. Charon is even similar in size to Uranus's moons Umbriel and Ariel. Charon's slow rotation means that there should be little flattening or tidal distortion if Charon is sufficiently massive to be in hydrostatic equilibrium. Any deviation from a perfect sphere is too small to have been detected by observations by the New Horizons mission. This is in contrast to Iapetus, a Saturnian moon similar in size to Charon but with a pronounced oblateness dating to early in its history. The lack of such oblateness in Charon could mean that it is currently in hydrostatic equilibrium, or simply that its orbit approached its current one early in its history, when it was still warm.
Based on mass updates from observations made by New Horizons the mass ratio of Charon to Pluto is 0.1218:1. This is much larger than the Moon to the Earth: 0.0123:1. Because of the high mass ratio, the barycenter is outside of the radius of Pluto, and the Pluto–Charon system has been referred to as a dwarf double planet. With four smaller satellites in orbit about the two larger worlds, the Pluto–Charon system has been considered in studies of the orbital stability of circumbinary planets.
Internal structure
Charon's volume and mass allow calculation of its density, , from which it can be determined that Charon is slightly less dense than Pluto and suggesting a composition of 55% rock to 45% ice (± 5%), whereas Pluto is about 70% rock. The difference is considerably lower than that of most suspected collisional satellites.
Following the New Horizons flyby, numerous discovered features on Charon's surface strongly indicated that Charon is differentiated, and may even have had a subsurface ocean early in its history. The past resurfacing observed on Charon's surface indicated that Charon's ancient subsurface ocean may have fed large-scale cryoeruptions on the surface, erasing many older features. As a result, two broad competing views on the nature of Charon's interior arose: the so-called hot start model, where Charon's formation is rapid and involves a violent impact with Pluto, and the cold start model, where Charon's formation is more gradual and involves a less violent impact with Pluto.
According to the hot start model, Charon accreted rapidly (within ~ years) from the circumplanetary disc, resulting from a highly-disruptive giant impact scenario. This rapid time scale prevents the heat from accretion from radiating away during the formation process, leading to the partial melting of Charon's outer layers. However, Charon's crust failed to reach a melt fraction where complete differentiation occurs, leading to the crust retaining part of its silicate content upon freezing. A liquid subsurface ocean forms during or soon after Charon's accretion and persists for approximately 2 billion years before freezing, possibly driving cryovolcanic resurfacing of Vulcan Planitia. Radiogenic heat from Charon's core could then melt a second subsurface ocean composed of a eutectic water-ammonia mixture before it too freezes, possibly driving the formation of Kubrick Mons and other similar features. These freezing cycles could increase Charon's size by >20 km, leading to the formation of the complex tectonic features observed in Serenity Chasma and Oz Terra.
In contrast, the cold start model argues that a large subsurface ocean early in Charon's history is not necessary to explain Charon's surface features, and instead proposes that Charon may have been homogeneous and more porous at formation. According to the cold start model, as Charon's interior begins to warm due to radiogenic heating and heating from serpentinization, a phase of contraction begins, largely driven by compaction in Charon's interior. Approximately 100-200 million years after formation, enough heat builds up to where a subsurface ocean melts, leading to rapid differentiation, further contraction, and the hydration of core rocks. Despite this melting, a pristine crust of amorphous water ice on Charon remains. After this period, differentiation continues, but the core can no longer absorb more water, and thus freezing at the base of Charon's mantle begins. This freezing drives a period of expansion until Charon's core becomes warm enough to begin compaction, starting a final period of contraction. Serenity Chasma may have formed from the expansion episode, whilst the final contraction episode may have given rise to the arcuate ridges observed in Mordor Macula.
Surface
Unlike Pluto's surface, which is composed of nitrogen and methane ices, Charon's surface appears to be dominated by the less volatile water ice.
In 2007, observations by the Gemini Observatory detected patches of ammonia hydrates and water crystals on the surface of Charon that suggested the presence of active cryogeysers and cryovolcanoes. The fact that the ice was still in crystalline form suggested it may have been deposited recently, as it was expected that solar radiation would have degraded it to an amorphous state after roughly thirty thousand years. However, following new data from the New Horizons flyby, no active cryovolcanoes or geysers were detected. Later research has also called into question the cryovolcanic origin for the crystalline water ice and ammonia features, with some researchers instead proposing that ammonia may be replenished passively from underground material.
Photometric mapping of Charon's surface shows a latitudinal trend in albedo, with a bright equatorial band and darker poles. The north polar region is dominated by a very large dark area informally dubbed "Mordor" by the New Horizons team. The favored explanation for this feature is that it is formed by condensation of gases that escaped from Pluto's atmosphere. In winter, the temperature is −258 °C, and these gases, which include nitrogen, carbon monoxide, and methane, condense into their solid forms; when these ices are subjected to solar radiation, they chemically react to form various reddish tholins. Later, when the area is again heated by the Sun as Charon's seasons change, the temperature at the pole rises to −213 °C, resulting in the volatiles sublimating and escaping Charon, leaving only the tholins behind. Over millions of years, the residual tholin builds up thick layers, obscuring the icy crust. In addition to Mordor, New Horizons found evidence of extensive past geology that suggests that Charon is probably differentiated; in particular, the southern hemisphere has fewer craters than the northern and is considerably less rugged, suggesting that a massive resurfacing event—perhaps prompted by the partial or complete freezing of an internal ocean—occurred at some point in the past and removed many of the earlier craters.
Charon has a system of extensive grabens and scarps, such as Serenity Chasma, which extend as an equatorial belt for at least . Argo Chasma potentially reaches as deep as , with cliffs that may rival Verona Rupes on Miranda for the title of the tallest cliff in the Solar System.
Hypothesized exosphere
In contrast to Pluto, Charon has no significant atmosphere. There has been speculation about an extremely thin exosphere surrounding the moon contributing to the formation of dark regions such as Mordor Macula. The strong seasons experienced by Pluto and Charon could provide brief periods of exosphere formation as methane sublimates on Charon, interspersed by centuries of dormancy.
Pluto does have a thin but significant atmosphere, which Charon's gravitation might pull toward Charon's surface. The gas, specifically nitrogen, is mostly caught in the combined center of gravity between the two bodies before reaching Charon, but any gas that does reach Charon is held closely against the surface. The gas is mostly made up of ions of nitrogen, but the amounts are negligible compared to the total of Pluto's atmosphere.
The many spectral signatures of ice formations on the surface of Charon have led some to believe that the ice formations could supply an atmosphere, but atmosphere supplying formations have not been confirmed yet. Many scientists theorize that these ice formations could be concealed out of direct sight, either in deep craters or beneath Charon's surface. Charon's relatively low gravity, due to its low mass, causes any atmosphere that might be present to rapidly escape the surface into space. Even through stellar occultation, which is used to probe the atmosphere of stellar bodies, scientists cannot confirm an existing atmosphere; this was tested in 1986 while attempting to perform stellar occultation testing on Pluto. Charon also acts as a protector for Pluto's atmosphere, blocking the solar wind that would normally collide with Pluto and damage its atmosphere. Since Charon blocks these solar winds, its own atmosphere is diminished, instead of Pluto's. This effect is also a potential explanation for Charon's lack of atmosphere; the solar winds remove gases faster than they can accumulate. It is still possible for Charon to have an atmosphere, as Pluto transfers some of its atmospheric gas to Charon, from where it tends to escape into space. Assuming Charon's density is 1.71 g/cm3, it would have a surface gravity of 0.6 of Pluto's. It also has a higher mean molecular weight than Pluto and a lower exobase surface temperature, so that the gases in its atmosphere would not escape as rapidly from Charon as they do from Pluto.
There has been significant proof of CO2 gas and H2O vapor on the surface of Charon, but these vapors are not sufficient for a viable atmosphere due to their low vapor pressures. Pluto's surface has abundant ice formations, but these are volatile, as they are made up of volatile substances like methane. These volatile ice structures cause a good deal of geological activity, keeping its atmosphere constant, while Charon's ice structures are mainly made up of water and carbon dioxide, much less volatile substances that can stay dormant and not affect the atmosphere much.
Observation and exploration
Since the first blurred images of the moon (1), images showing Pluto and Charon resolved into separate disks were taken for the first time by the Hubble Space Telescope in the 1990s (2). The telescope was responsible for the best, yet low-quality images of the moon. In 1994, the clearest picture of the Pluto–Charon system showed two distinct and well-defined disks (3). The image was taken by Hubble's Faint Object Camera (FOC) when the system was 4.4 billion kilometers (2.6 billion miles) away from Earth Later, the development of adaptive optics made it possible to resolve Pluto and Charon into separate disks using ground-based telescopes. Although ground-based observation is very challenging, a group of amateur astronomers in Italy used a 14-inch telescope in 2008 to successfully resolve Charon in an image of Pluto.
In June 2015, the New Horizons spacecraft captured consecutive images of the Pluto–Charon system as it approached it. The images were put together in an animation. It was the best image of Charon to that date (4). In July 2015, the New Horizons spacecraft made its closest approach to the Pluto system. It is the only spacecraft to date to have visited and studied Charon. Charon's discoverer James Christy and the children of Clyde Tombaugh were guests at the Johns Hopkins Applied Physics Laboratory during the New Horizons closest approach.
Classification
The center of mass (barycenter) of the Pluto–Charon system lies outside either body. Because neither object truly orbits the other, and Charon has 12.2% of the mass of Pluto, it has been argued that Charon should be considered to be part of a binary planet with Pluto. The International Astronomical Union (IAU) states that Charon is a satellite of Pluto, but the idea that Charon might be classified as a dwarf planet in its own right may be considered at a later date.
In a draft proposal for the 2006 redefinition of the term, the IAU proposed that a planet is defined as a body that orbits the Sun that is large enough for gravitational forces to render the object (nearly) spherical. Under this proposal, Charon would have been classified as a planet, because the draft explicitly defined a planetary satellite as one in which the barycenter lies within the major body. In the final definition, Pluto was reclassified as a dwarf planet, but a formal definition of a planetary satellite was not decided upon. Charon is not in the list of dwarf planets currently recognized by the IAU. Had the draft proposal been accepted, even the Moon would hypothetically be classified as a planet in billions of years when the tidal acceleration that is gradually moving the Moon away from Earth takes it far enough away that the center of mass of the system no longer lies within Earth.
The other moons of Pluto Nix, Hydra, Kerberos, and Styx orbit the same barycenter but they are not large enough to be spherical and they are simply considered to be satellites of Pluto (or of Pluto–Charon).
| Physical sciences | Solar System | Astronomy |
52866 | https://en.wikipedia.org/wiki/Cooking%20weights%20and%20measures | Cooking weights and measures | In recipes, quantities of ingredients may be specified by mass (commonly called weight), by volume, or by count.
For most of history, most cookbooks did not specify quantities precisely, instead talking of "a nice leg of spring lamb", a "cupful" of lentils, a piece of butter "the size of a small apricot", and "sufficient" salt. Informal measurements such as a "pinch", a "drop", or a "hint" (soupçon) continue to be used from time to time. In the US, Fannie Farmer introduced the more exact specification of quantities by volume in her 1896 Boston Cooking-School Cook Book.
Today, most of the world prefers metric measurement by weight, though the preference for volume measurements continues among home cooks in the United States and the rest of North America. Different ingredients are measured in different ways:
Liquid ingredients are generally measured by volume worldwide.
Dry bulk ingredients, such as sugar and flour, are measured by weight in most of the world ("250 g flour"), and by volume in North America (" cup flour"). Small quantities of salt and spices are generally measured by volume worldwide, as few households have sufficiently precise balances to measure by weight.
In most countries, meat is described by weight or count: "a 2 kilogram chicken"; "four lamb chops".
Eggs are usually specified by count. Vegetables are usually specified by weight or occasionally by count, despite the inherent imprecision of counts given the variability in the size of vegetables.
Metric measures
In most of the world, recipes use the metric system of units—litres (L) and millilitres (mL), grams (g) and kilograms (kg), and degrees Celsius (°C). The official spelling litre is used in most English-speaking nations; the notable exception is the United States where the spelling liter is preferred.
The United States measures weight in pounds (avoirdupois), while recipes in the UK tend to include both imperial and metric measures, following the advice of the Guild of Food Writers. The United States also uses volume measures based on cooking utensils and pre-metric measures. The actual values frequently deviate from the utensils on which they were based, and there is little consistency from one country to another.
In South Australia, a "pint" of beer is traditionally 425 mL, while most other states have metricated this value to 570 mL.
‡ In Canada, a cup was historically 8 imperial fluid ounces (227 mL) but could also refer to 10 imperial fl oz (284 mL), as in Britain, and even a metric cup of 250 mL. Serving sizes on nutrition labelling on food packages in Canada employ the metric cup of 250 mL, with nutrition labelling in the US using a cup of 240 mL, based on the US customary cup.
* In the UK, teaspoons and tablespoons are formally and of an imperial pint (3·55 mL and 14·21 mL), respectively. In Canada, a teaspoon is historically imperial fluid ounce (4.74 mL) and a tablespoon is imperial fl oz (14.21 mL). In both Britain and Canada, cooking utensils commonly come in 5 mL for teaspoons and 15 mL for tablespoons, hence why it is labelled as that on the chart.
The volume measures here are for comparison only. See below for the definition of Gallon for more details.
In addition, the "cook's cup" above is not the same as a "coffee cup", which can vary anywhere from , or even smaller for espresso.
In Australia, since 1970, metric utensil units have been standardized by law, and imperial measures no longer have legal status. However, it is wise to measure the actual volume of the utensil measures, particularly the 'Australian tablespoon' (see above), since many are imported from other countries with different values. Dessertspoons are standardized as part of the metric system at 10 mL, though they are not normally used in contemporary recipes. Australia is the only metricated country with a metric tablespoon of 20 mL, unlike other countries that metricated, which have a 15 mL metric tablespoon.
In Europe, older recipes frequently refer to "pounds" (e.g. in German, in Dutch, in French). In each case, the unit refers to 500 g, about 10% more than an avoirdupois pound (454 g). Dutch recipes may also use the , which is 100 g.
Weight of liquids
With the advent of accurate electronic scales, it has become more common to weigh liquids for use in recipes, avoiding the need for accurate volumetric utensils. The most common liquids used in cooking are water and milk, milk having approximately the same density as water.
1 mL of water weighs 1 gram so a recipe calling for 300 mL (≈ Imperial Pint) of water can simply be substituted with 300 g (≈ 10 oz.) of water.
1 fluid ounce of water weighs approximately 1 ounce so a recipe calling for a UK pint (20 fl oz) of water can be substituted with 20 oz of water.
More accurate measurements become important in the large volumes used in commercial food production. Also, a home cook can use greater precision at times. Water at may be volumetrically measured then weighed to determine an unknown measuring-utensil volume without the need for a water-density adjustment.
United States measures
The US uses pounds and ounces (avoirdupois) for weight, and US customary units for volume. For measures used in cookbooks published in other nations navigate to the appropriate regional section in Traditional measurement systems.
Measures are classified as either dry measures or fluid measures. Some of the fluid and dry measures have similar names, but the actual measured volume is quite different. A recipe will generally specify which measurement is required. U.S. recipes are commonly in terms of fluid measures, even for dry ingredients. Most of these units derive from earlier English units, as applied to the U.S. gallon. Typically they follow a pattern of binary submultiples, where each larger measure consists of two units of the next-smallest measure. An exception is with the commonly used teaspoon as one-third of a tablespoon.
Binary submultiples are fractional parts obtained by successively dividing by the number 2. Thus, one-half, one-fourth, one-eighth, one-sixteenth, and so on, are binary submultiples. The system can be traced back to the measuring systems of the Hindus and the ancient Egyptians, who subdivided the (about 4.8 litres) into parts of , ,, , , and (1 , or mouthful, or about 14.5 ml), and the similarly down to (1 ) using hieratic notation, as early as the Fifth Dynasty of Egypt, 2494 to 2345 BC, thus making the "English doubling system" at least 4300 years old.
From units and tools of convenience, most of the system's history could have values vary widely, and was not until recent centuries that standardization began to take shape. Moreso, the overlap with other systems like the apothecaries' system, and giving each division a unique and often variable by context, person, place, and time name instead of a systematic one, can make the system seem confusing for those not accustomed to it. However, other than the names themselves, the regular ratios make the actual measurements straightforward; and in many cases names have been deprecated in favor of fractionally denominated amounts of a few core units (such as taking gallons, cups, and teaspoons to their nearest quarters without names: nixing pottle; gill and wineglass; dram (as a culinary unit), coffeespoon, and saltspoon; respectively), or are limited to the specific or esoteric.
It is still a legal basis for measures in many states, such as Massachusetts, which mandates that "Glass bottles or jars used for the sale of milk or cream to the consumer shall be of the capacity of one gallon, a multiple of the gallon, or a binary submultiple of the gallon."
Metric equivalents are based upon one of two nearly equivalent systems. In the standard system the conversion is that 1 gallon = 231 cubic inches and 1 inch = 2.54 cm, which makes a gallon = 3785.411784 millilitres exactly. For nutritional labeling on food packages in the US, the teaspoon is defined as exactly 5 ml, giving 1 gallon = 3840 ml exactly. This chart uses the former.
Suffixed asterisks on some of the "tsp" units in the "Defined" column above indicate that those teaspoon units are defined as fl oz (4 fl dram), the old 4 tsp = 1 tbsp amount, instead of fl oz. This definition fits with "barkeepers' teaspoon", and is used in many cocktail recipe books; generally the subdivisions are not so explicitly defined nor named below tsp in general culinary. This can be verified by comparing the associated values in the "fl oz" column. All other "tsp" units in the "Defined" column are indeed defined as fl oz, the current 3 tsp = 1 tbsp amount.
* Discrepancies due to size, generally disregarded as at the scale it becomes a factor, the person generally is using the next size up measuring cup (i.e.: fl oz is likely to be straight measured in an ounce cup and not as 9 (vs 12) teaspoons)
‡ Rare if not nonexistent in use by name rather than as fraction of a different unit.
† The fluid scruple has been properly defined on its own in the apothecaries' system as fl oz, fluid dram, or = 20 minims (≈ 1.23223 ml), and also tsp. Mind that scruples and drams were pharmaceutical and intended to be specific and precise, whereas cooking measures tended to use what was on hand and/or actually used to consume what was being prepared, and not intended to be as formally scientific in its degree of precision.
The saltspoon most likely combined into the scruple over time, as a consequence of home cooks approximating standard measures with what they had at hand, much as the teaspoon was roughly "close enough" for a kitchen approximation to a fluid dram (= 60 minims), but not equal to the fl dr (80 minims) value it actually is; especially with the variability of the method of measuring itself. Not of insignificance is the natural habit of customary measures to use a 2n dividing scheme regardless of exact definitions, and this pattern is seen even with metric measuring spoons
Confusion comes about from teaspoon continuing to be called a "dram" in vernacular, despite the sizes of actual spoons creeping up quietly over time, such that of a tsp (tsp > fl dr) would in fact become congruent with current fl dr values for the scruple and saltspoon; in other words, the terminology not keeping pace with the definition. At the small scales involved this is negligible (i.e.: math can convert down to tsp ×10−9, but to what degree can it practically be meted); however can cause problems when accuracy is required such as medicines: "In almost all cases the modern teacups, tablespoons, dessertspoons, and teaspoons, after careful test by the author, were found to average 25 percent greater capacity than the theoretical quantities given above, and thus the use of accurately graduated medicine glasses, which may be had now at a trifling cost, should be insisted upon."
* Discrepancies due to size, generally disregarded as at the scale it becomes a factor, the person generally is using the next size up measuring cup (i.e.: fl oz is likely to be straight measured in an ounce cup and not as 9 (vs 12) teaspoons)
In domestic cooking, bulk solids, notably flour and sugar, are measured by volume, often cups, though they are sold by weight at retail. Weight measures are used for meat. Butter may be measured by either weight ( lb) or volume (3 tbsp) or a combination of weight and volume ( lb plus 3 tbsp); it is sold by weight but in packages marked to facilitate common divisions by eye. (As a sub-packaged unit, a stick of butter, at lb [113 g], is a de facto measure in the US.) Some recipes may specify butter amounts called a pat (1 - 1.5 tsp) or a knob (2 tbsp).
Cookbooks in Canada use the same system, although pints and gallons would be taken as their Imperial quantities unless specified otherwise. Following the adoption of the metric system, recipes in Canada are frequently published with metric conversions.
Approximate units
There are a variety of approximate units of measures, which are frequently undefined by any official source, or which have had conflicting definitions over time, yet which are commonly used. The measurement units that are most commonly understood to be approximate are the drop, smidgen, pinch, and dash, yet nearly all of the traditional cooking measurement units lack statutory definitions, or even any definition by any organization authorized to set standards in the U.S. For example, of the table above, only the fluid ounce, pint, quart, and gallon are officially defined by the NIST. All of the others appear only in conversion guides lacking statutory authority, or in now-obsolete publications of the U.S. Pharmacopeial Convention, or USP—essentially, the Apothecaries' system—which still has authority to define certain drug and supplement standards. The USP has long-since abandoned Apothecaries' measurements, and even now recommends against using teaspoons to measure doses of medicine.
British (Imperial) measures
Note that measurements in this section are in imperial units.
British imperial measures distinguish between weight and volume.
Weight is measured in ounces and pounds (avoirdupois) as in the U.S.
Volume is measured in imperial gallons, quarts, pints, fluid ounces, fluid drachms, and minims. The imperial gallon was originally defined as of water in 1824, and refined as exactly 4.54609 litres in 1985.
Traditionally, when describing volumes, recipes commonly give measurements in the following units:
Tumbler (10 fluid ounces; named after a typical drinknig glass)
Breakfast cup (8 fluid ounces; named after a cup for drinking tea or coffee while eating breakfast)
Cup (6 fluid ounces; named after an everyday drinking cup)
Teacup (5 fluid ounces; named after a typical teacup)
Coffee cup ( fluid ounces; named after a small cup for serving after‑dinner coffee)
Wine glass (2 fluid ounces; named after a small glass for serving liquor)
If the recipe is one that has been handed down in a family and gives measurements in ‘cups’, it is just as likely to refer to someone's favourite kitchen cup as to the said unit that is 6 fluid ounces.
All six units are the traditional British equivalents of the US customary cup and the metric cup, used in situations where a US cook would use the US customary cup and a cook using metric units the metric cup. The breakfast cup is the most similar in size to the US customary cup and the metric cup. Which of these six units is used depends on the quantity or volume of the ingredient: there is division of labour between these six units, like the tablespoon and the teaspoon. British cookery books and recipes, especially those from the days before the UK's partial metrication, commonly use two or more of the aforesaid units simultaneously: for example, the same recipe may call for a ‘tumblerful’ of one ingredient and a ‘wineglassful’ of another one; or a ‘breakfastcupful’ or ‘cupful’ of one ingredient, a ‘teacupful’ of a second one, and a ‘coffeecupful’ of a third one. Unlike the US customary cup and the metric cup, a tumbler, a breakfast cup, a cup, a teacup, a coffee cup, and a wine glass are not measuring cups: they are simply everyday drinking vessels commonly found in British households and typically having the respective aforementioned capacities; due to long-term and widespread use, they have been transformed into measurement units for cooking. There is not a British imperial unit–based culinary measuring cup.
For smaller amounts, British recipes traditionally give measurements in the following units:
Tablespoon (4 fluid drachms or fluid ounce)
Dessert spoon ( tablespoon: the equivalence of 2 fluid drachms or fluid ounce)
Teaspoon ( dessert spoon or tablespoon: the equivalence of 1 fluid drachm or fluid ounce)
Salt spoon ( teaspoon: the equivalence of 30 minims, fluid drachm, or fluid ounce)
For even smaller amounts, the following units are used:
Pinch ( salt spoon or teaspoon: an amount of space that can accommodate 15 minims ( fluid drachm or fluid ounce) of liquid), if it is a dry ingredient
Drop (1 minim, fluid drachm, or fluid ounce), if it is a liquid
American cooks using British recipes, and vice versa, need to be careful with pints and fluid ounces.
A US pint (16 US fluid ounces) is about 16·65 UK fluid ounces or 473 mL, while a UK pint is 20 UK fluid ounces (about 19·21 US fluid ounces or 568 mL): a UK pint is, therefore, about 20% larger than a US pint.
A US fluid ounce is of a US pint (about 1·04 UK fluid ounces or 29.6 mL); a UK fluid ounce is of a UK pint (about 0·96 US fluid ounce or 28.4 mL).
On a larger scale, perhaps for institutional cookery, a UK gallon is 8 UK pints (160 UK fluid ounces; about 1·2 US gallons or 4.546 litres), whereas the US gallon is 8 US pints (128 US fluid ounces; about 0·83 UK gallon or 3.785 litres).
The metric system was officially adopted in the UK, for most purposes, in the 20th century and both imperial and metric are taught in schools and used in books. It is now mandatory for the sale of food to also show metric. However, it is not uncommon to purchase goods which are measured and labeled in metric, but the actual measure is rounded to the equivalent imperial measure (i.e., milk labeled as 568 mL / 1 pint). In September 2007, the EU with Directive 2007/45/EC deregulated prescribed metric packaging of most products, leaving only wines and liqueurs subject to prescribed EU-wide pre-packaging legislation; the law relating to labelling of products remaining unchanged.
Special instructions
Volume measures of compressible ingredients have a substantial measurement uncertainty, in the case of flour of about 20%. Some volume-based recipes, therefore, attempt to improve the reproducibility by including additional instructions for measuring the correct amount of an ingredient. For example, a recipe might call for "1 cup brown sugar, firmly packed", or "2 heaping cups flour". A few of the more common special measuring methods:
Firmly packed
With a spatula, a spoon, or by hand, the ingredient is pressed as tightly as possible into the measuring device.
Lightly packed
The ingredient is pressed lightly into the measuring device, only tightly enough to ensure no air pockets.
Even / level
A precise measure of an ingredient, discarding all of the ingredient that rises above the rim of the measuring device. Sweeping across the top of the measure with the back of a straight knife or the blade of a spatula is a common leveling method.
Rounded
Allowing a measure of an ingredient to pile up above the rim of the measuring device naturally, into a soft, rounded shape.
Heaping / heaped
The maximum amount of an ingredient which will stay on the measuring device.
Sifted
This instruction may be seen in two different ways, with two different meanings: before the ingredient, as "1 cup sifted flour", indicates the ingredient should be sifted into the measuring device (and normally leveled), while after the ingredient, as "1 cup flour, sifted", denotes the sifting should occur after measurement.
Such special instructions are unnecessary in weight-based recipes.
| Physical sciences | Measurement systems | Basics and measurement |
52957 | https://en.wikipedia.org/wiki/Bell%27s%20palsy | Bell's palsy | Bell's palsy is a type of facial paralysis that results in a temporary inability to control the facial muscles on the affected side of the face. In most cases, the weakness is temporary and significantly improves over weeks. Symptoms can vary from mild to severe. They may include muscle twitching, weakness, or total loss of the ability to move one or, in rare cases, both sides of the face. Other symptoms include drooping of the eyebrow, a change in taste, and pain around the ear. Typically symptoms come on over 48 hours. Bell's palsy can trigger an increased sensitivity to sound known as hyperacusis.
The cause of Bell's palsy is unknown and it can occur at any age. Risk factors include diabetes, a recent upper respiratory tract infection, and pregnancy. It results from a dysfunction of cranial nerve VII (the facial nerve). Many believe that this is due to a viral infection that results in swelling. Diagnosis is based on a person's appearance and ruling out other possible causes. Other conditions that can cause facial weakness include brain tumor, stroke, Ramsay Hunt syndrome type 2, myasthenia gravis, and Lyme disease.
The condition normally gets better by itself, with most achieving normal or near-normal function. Corticosteroids have been found to improve outcomes, while antiviral medications may be of a small additional benefit. The eye should be protected from drying up with the use of eye drops or an eyepatch. Surgery is generally not recommended. Often signs of improvement begin within 14 days, with complete recovery within six months. A few may not recover completely or have a recurrence of symptoms.
Bell's palsy is the most common cause of one-sided facial nerve paralysis (70%). It occurs in 1 to 4 per 10,000 people per year. About 1.5% of people are affected at some point in their lives. It most commonly occurs in people between ages 15 and 60. Males and females are affected equally. It is named after Scottish surgeon Charles Bell (1774–1842), who first described the connection of the facial nerve to the condition.
Although defined as a mononeuritis (involving only one nerve), people diagnosed with Bell's palsy may have "myriad neurological symptoms", including "facial tingling, moderate or severe headache/neck pain, memory problems, balance problems, ipsilateral limb paresthesias, ipsilateral limb weakness, and a sense of clumsiness" that are "unexplained by facial nerve dysfunction".
Signs and symptoms
Bell's palsy is characterized by a one-sided facial droop that comes on within 72 hours. In rare cases (<1%), it can occur on both sides resulting in total facial paralysis.
The facial nerve controls many functions, such as blinking and closing the eyes, smiling, frowning, lacrimation, salivation, flaring nostrils and raising eyebrows. It also carries taste sensations from the anterior two thirds of the tongue, through the chorda tympani nerve (a branch of the facial nerve). Because of this, people with Bell's palsy may present with loss of taste sensation in the anterior two thirds of the tongue on the affected side.
Although the facial nerve innervates the stapedius muscle of the middle ear (through the tympanic branch), sound sensitivity, causing normal sounds to be perceived as very loud (hyperacusis), and dysacusis is possible but hardly ever clinically evident.
Cause
The cause of Bell's palsy is unknown. Risk factors include diabetes, a recent upper respiratory tract infection, and pregnancy.
Some viruses are thought to establish a persistent (or latent) infection without symptoms, e.g., the varicella zoster virus and the Epstein–Barr virus, both of the herpes family. Reactivation of an existing (dormant) viral infection has been suggested as a cause of acute Bell's palsy. As the facial nerve swells and becomes inflamed in reaction to the infection, it causes pressure within the Fallopian canal, resulting in the restriction of blood and oxygen to the nerve cells. Other viruses and bacteria that have been linked to the development of Bell's palsy include HIV, sarcoidosis and Lyme disease. This new activation could be triggered by trauma, environmental factors, and metabolic or emotional disorders.
Familial inheritance has been found in 4–14% of cases. There may also be an association with migraines.
In December 2020, the U.S. FDA recommended that recipients of the Pfizer and Moderna COVID-19 vaccines should be monitored for symptoms of Bell's palsy after several cases were reported among clinical trial participants, though the data were not sufficient to determine a causal link.
Genetics
A meta-analysis of genome-wide association study (GWAS) identified the first unequivocal association with Bell's palsy.
Pathophysiology
Bell's palsy is the result of a malfunction of the facial nerve (cranial nerve VII), which controls the muscles of the face. Facial palsy is typified by an inability to move the muscles of facial expression. The paralysis is of the infranuclear/lower motor neuron type.
It is thought that as a result of inflammation of the facial nerve, pressure is produced on the nerve where it exits the skull within its bony canal (the stylomastoid foramen), blocking the transmission of neural signals or damaging the nerve. Patients with facial palsy for which an underlying cause can be found are not considered to have Bell's palsy per se. Possible causes of facial paralysis include tumor, meningitis, stroke, diabetes mellitus, head trauma and inflammatory diseases of the cranial nerves (sarcoidosis, brucellosis, etc.). In these conditions, the neurologic findings are rarely restricted to the facial nerve. Babies can be born with facial palsy. In a few cases, bilateral facial palsy has been associated with acute HIV infection.
In some research, the herpes simplex virus type 1 (HSV-1) has been identified in a majority of cases diagnosed as Bell's palsy through endoneurial fluid sampling. Other research, however, identified, out of a total of 176 cases diagnosed as Bell's palsy, HSV-1 in 31 cases (18%) and herpes zoster in 45 cases (26%).
In addition, HSV-1 infection is associated with demyelination of nerves. This nerve damage mechanism is different from the above-mentioned—that edema, swelling, and compression of the nerve in the narrow bone canal are responsible for nerve damage. Demyelination may not even be directly caused by the virus but by an unknown immune response.
Diagnosis
Bell's palsy is a diagnosis of exclusion, meaning it is diagnosed by the elimination of other reasonable possibilities. By definition, no specific cause can be determined. There are no routine lab or imaging tests required to make the diagnosis. The degree of nerve damage can be assessed using the House-Brackmann score.
One study found that 45% of patients are not referred to a specialist, which suggests that Bell's palsy is considered by physicians to be a straightforward diagnosis that is easy to manage.
Other conditions that can cause similar symptoms include herpes zoster, Lyme disease, sarcoidosis, stroke, and brain tumors.
Differential diagnosis
Once the facial paralysis sets in, many people may mistake it as a symptom of a stroke; however, there are a few subtle differences. A stroke will usually cause a few additional symptoms, such as numbness or weakness in the arms and legs. And unlike Bell's palsy, a stroke will usually let patients control the upper part of their faces. A person with a stroke will usually have some wrinkling on their forehead.
In areas where Lyme disease is common, it accounts for about 25% of cases of facial palsy. In the U.S., Lyme is most common in the New England and Mid-Atlantic states and parts of Wisconsin and Minnesota. The first sign of about 80% of Lyme infections, typically one or two weeks after a tick bite, is usually an expanding rash that may be accompanied by headaches, body aches, fatigue, or fever. In up to 10–15% of Lyme infections, facial palsy appears several weeks later, and may be the first sign of infection that is noticed as the Lyme rash typically does not itch and is not painful. The likelihood that the facial palsy is caused by Lyme disease should be estimated, based on the recent history of outdoor activities in likely tick habitats during warmer months, a recent history of rash or symptoms such as headache and fever, and whether the palsy affects both sides of the face (much more common in Lyme than in Bell's palsy). If that likelihood is more than negligible, a serological test for Lyme disease should be performed, and if it exceeds 10%, empiric therapy with antibiotics should be initiated, without corticosteroids, and reevaluated upon completion of laboratory tests for Lyme disease. Corticosteroids have been found to harm outcomes for facial palsy caused by Lyme disease.
One disease that may be difficult to exclude in the differential diagnosis is the involvement of the facial nerve in infections with the herpes zoster virus. The major differences in this condition are the presence of small blisters, or vesicles, on the external ear, significant pain in the jaw, ear, face, and/or neck, and hearing disturbances, but these findings may occasionally be lacking (zoster sine herpete). Reactivation of existing herpes zoster infection leading to facial paralysis in a Bell's palsy type pattern is known as Ramsay Hunt syndrome type 2. The prognosis for Bell's palsy patients is generally much better than for Ramsay Hunt syndrome type 2 patients.
Treatment
Steroids are effective at improving recovery in Bell's palsy while antivirals have not. In those who are unable to close their eyes, eye-protective measures are required. Management during pregnancy is similar to management in the non-pregnant.
Steroids
Corticosteroids such as prednisone improve recovery at 6 months and are thus recommended. Early treatment (within 3 days after the onset) is necessary for benefit with a 14% greater probability of recovery. There is some debate regarding the optimal dosing strategy which is generally physician dependent.
Antivirals
One review found that antivirals (such as aciclovir) are ineffective in improving recovery from Bell's palsy beyond steroids alone in mild to moderate disease. Another review found a benefit when combined with corticosteroids but stated the evidence was not very good to support this conclusion.
In severe disease, it is also unclear. One 2015 review found no effect regardless of severity. Another review found a small benefit when added to steroids.
They are commonly prescribed due to a theoretical link between Bell's palsy and the herpes simplex and varicella zoster virus. There is still the possibility that they might result in a benefit less than 7% as this has not been ruled out.
Eye protection
When Bell's palsy affects the blink reflex and stops the eye from closing completely, frequent use of tear-like eye drops or eye ointments is recommended during the day, and protecting the eyes with patches or taping them shut is recommended for sleep and rest periods.
Physiotherapy
Physiotherapy can be beneficial to some individuals with Bell's palsy as it helps to maintain muscle tone of the affected facial muscles and stimulate the facial nerve. It is important that muscle re-education exercises and soft tissue techniques be implemented before recovery to help prevent permanent contractures of the paralyzed facial muscles. To reduce pain, heat can be applied to the affected side of the face. There is no high-quality evidence to support the role of electrical stimulation for Bell's palsy.
Surgery
Surgery may be able to improve outcomes in facial nerve palsy that has not recovered. A number of different techniques exist. Smile surgery or smile reconstruction is a surgical procedure that may restore the smile for people with facial nerve paralysis. Adverse effects include hearing loss which occurs in 3–15% of people. A Cochrane review (updated in 2021), after reviewing applicable randomized and quasi-randomized controlled trials was unable to determine if early surgery is beneficial or harmful. As of 2007 the American Academy of Neurology did not recommend surgical decompression.
Alternative medicine
The efficacy of acupuncture remains unknown because the available studies are of low quality (poor primary study design or inadequate reporting practices). There is very tentative evidence for hyperbaric oxygen therapy in severe disease.
Prognosis
Most people with Bell's palsy start to regain normal facial function within three weeks—even those who do not receive treatment. In a 1982 study, when no treatment was available, of 1,011 patients, 85% showed first signs of recovery within three weeks after onset. For the other 15%, recovery occurred 3–6 months later.
After a follow-up of at least one year or until restoration, complete recovery had occurred in more than two-thirds (71%) of all patients. Recovery was judged moderate in 12% and poor in only 4% of patients. Another study found that incomplete palsies disappear entirely, nearly always in one month. The patients who regain movement within the first two weeks nearly always remit entirely. When remission does not occur until the third week or later, a significantly greater part of the patients develop sequelae. A third study found a better prognosis for young patients, aged below 10 years old, while the patients over 61 years old presented a worse prognosis.
Major possible complications of the condition are chronic loss of taste (ageusia), chronic facial spasm, facial pain, and corneal infections. Another complication can occur in case of incomplete or erroneous regeneration of the damaged facial nerve. The nerve can be thought of as a bundle of smaller individual nerve connections that branch out to their proper destinations. During regrowth, nerves are generally able to track the original path to the right destination—but some nerves may sidetrack leading to a condition known as synkinesis. For instance, the regrowth of nerves controlling muscles attached to the eye may sidetrack and also regrow connections reaching the muscles of the mouth. In this way, the movement of one also affects the other. For example, when the person closes the eye, the corner of the mouth lifts involuntarily.
Around 9% of people have some sort of ongoing problems after Bell's palsy, typically the synkinesis already discussed, or spasm, contracture, tinnitus, or hearing loss during facial movement or crocodile-tear syndrome. This is also called gustatolacrimal reflex or Bogorad's syndrome and results in shedding tears while eating. This is thought to be due to faulty regeneration of the facial nerve, a branch of which controls the lacrimal and salivary glands. Gustatorial sweating can also occur.
Epidemiology
The number of new cases of Bell's palsy ranges from about one to four cases per 10,000 population per year. The rate increases with age. Bell's palsy affects about 40,000 people in the United States every year. It affects approximately 1 person in 65 during a lifetime.
A range of annual incidence rates have been reported in the literature: 15, 24, and 25–53 (all rates per 100,000 population per year). Bell's palsy is not a reportable disease, and there are no established registries for people with this diagnosis, which complicates precise estimation.
Frequency
About 40,000 people are affected by Bell's Palsy in the United States every year. It can affect anyone of any gender and age, but its incidence seems to be highest in those in the 15- to 45-year-old age group.
History
The Persian physician Muhammad ibn Zakariya al-Razi (865–925) detailed the first known description of peripheral and central facial palsy.
Cornelis Stalpart van der Wiel (1620–1702) in 1683 gave an account of Bell's palsy and credited the Persian physician Ibn Sina (980–1037) for describing this condition before him. James Douglas (1675–1742) and (1761–1836) also described it.
Scottish neurophysiologist Sir Charles Bell read his paper to the Royal Society of London on July 12, 1821, describing the role of the facial nerve. He became the first to detail the neuroanatomical basis of facial paralysis. Since then, idiopathic peripheral facial paralysis has been referred to as Bell's palsy, named after him.
A notable person with Bell's palsy is former Prime Minister of Canada Jean Chrétien. During the 1993 Canadian federal election, Chrétien's first as leader of the Liberal Party of Canada, the opposition Progressive Conservative Party of Canada ran an attack ad in which voice actors criticized him over images that seemed to highlight his abnormal facial expressions. The ad was interpreted as an attack on Chrétien's physical appearance and garnered widespread anger among the public, while Chrétien used the ad to make himself more sympathetic to voters. The ad had the adverse effect of increasing Chrétien's lead in the polls and the subsequent backlash clinched the election for the Liberals, which the party won in a landslide.
On April 25, 2024, Joel Hans Embiid of the Philadelphia 76ers scored 50 points in Game 3 of the 76ers' first-round playoff series versus the New York Knicks while suffering from a mild case of Bell's Palsy.
| Biology and health sciences | Disabilities | Health |
52963 | https://en.wikipedia.org/wiki/Obstetrics%20and%20gynaecology | Obstetrics and gynaecology | Obstetrics and gynaecology (also spelled as obstetrics and gynecology; abbreviated as Obst and Gynae, O&G, OB-GYN and OB/GYN) is the medical specialty that encompasses the two subspecialties of obstetrics (covering pregnancy, childbirth, and the postpartum period) and gynaecology (covering the health of the female reproductive system – vagina, uterus, ovaries, and breasts). The specialization is an important part of care for women's health.
Postgraduate training programs for both fields are usually combined, preparing the practising obstetrician-gynecologist to be adept both at the care of female reproductive organs' health and at the management of pregnancy, although many doctors go on to develop subspecialty interests in one field or the other.
Scope
United States
According to the American Board of Obstetrics and Gynecology (ABOG), which is responsible for issuing OB-GYN certifications in the United States, the first step to OB-GYN certification is completing medical school to receive an MD or DO degree. From there doctors must complete a four-year OB-GYN residency program approved by the Accreditation Council for Graduate Medical Education (ACGME). For the 2021 Electronic Residency Application Service (ERAS) match, there were 277 OB-GYN residency programs accepting applicants.
In their fourth year of residency, with an affidavit from their director to confirm program completion, OB-GYN residents can choose whether to begin the board certification process by applying to take the ABOG Qualifying Exam, which is a written test. If residents pass the Qualifying Exam, demonstrating they possess the knowledge and skills to potentially become certified OB-GYNs, they are then eligible to sit for the oral Certification Exam. Prior to the Certification Exam, residents must also gather a list of patient cases they've worked on throughout their residency in order to demonstrate their competence and experience in OB-GYN patient care.
Residents then sit for the three-hour oral exam at ABOG's test center, and if they pass the exam they become "board certified" OB-GYNs. Since 2013 at least 82% of all Certifying Exam examinees have passed.
This adds up to 11–14 years of education and practical experience. The first 7–9 years are general medical training.
Experienced OB-GYN professionals can seek certifications in sub-specialty areas, including maternal and fetal medicine. See Fellowship (medicine).
United Kingdom
All doctors must first complete medical school and obtain a MBBS or equivalent certification. This portion typically takes five years. Following this, they are eligible for provisional registration with the General Medical Council. They then must complete a two years of foundation training. After the first year of training is complete, trainees are eligible for full registration with the General Medical Council. After the foundation training is complete applicants take the Part 1 MRCOG examination administered by the Royal College of Obstetricians and Gynaecologists. There are an additional seven years of training after this, and two more exams (Part 2 and Part 3 MRCOG exams) which adds up to nine years total minimum in training, although some trainees may take longer.
Subspecialties
Examples of subspecialty training available to physicians in the US are:
Maternal-fetal medicine: an obstetrical subspecialty, sometimes referred to as perinatology, that focuses on the medical and surgical management of high-risk pregnancies and surgery on the fetus with the goal of reducing morbidity and mortality.
Reproductive endocrinology and infertility: a subspecialty that focuses on the biological causes and interventional treatment of infertility
Gynecological oncology: a gynaecologic subspecialty focusing on the medical and surgical treatment of women with cancers of the reproductive organs
Female pelvic medicine and reconstructive surgery: a gynaecologic subspecialty focusing on the diagnosis and surgical treatment of women with urinary incontinence and prolapse of the pelvic organs. Sometimes referred to by laypersons as "female urology"
Advanced laparoscopic surgery
Family planning: a gynaecologic subspecialty offering training in contraception and pregnancy termination (abortion)
Pediatric and adolescent gynaecology
Menopausal and geriatric gynaecology
Of these, only the first four are truly recognized sub-specialties by the Accreditation Council for Graduate Medical Education (ACGME) and the American Board of Obstetrics and Gynecology (ABOG). The other subspecialties are recognized as informal concentrations of practice. To be recognized as a board-certified subspecialist by the American Board of Obstetrics and Gynecology or the American Osteopathic Board of Obstetrics and Gynecology, a practitioner must have completed an ACGME or AOA-accredited residency and obtained a Certificate of Added Qualifications (CAQ) which requires an additional standardized examination.
Additionally, physicians of other specialties may become trained in Advanced Life Support in Obstetrics (ALSO), a short certification that equips them to better manage emergent OB/GYN situations.
Common procedures
There are many procedures that can be provided to people by OB/GYNs. Some procedures may include:
Colposcopy: If the results of a cervical cancer screening test, such as Pap smear or HPV test, are abnormal this more thorough examination of the cervix and vaginal tissues may be needed.
Loop electrical excision procedure (LEEP): a procedure to quickly remove abnormal vaginal tissue within the cervix. A local anesthetic and a solution to enhance the points of removal visually is administered during the process. There is a chance of experiencing watery, pinkish discharge, brownish discharge, and mild cramping.
Endometrial biopsy: a procedure that collects a tissue sample from the endometrium lining of the uterus. The sample is tested and checked under a microscope for abnormals cells or indicators of cancer.
IUD insertion: an intrauterine device that is T-shaped and is placed in the uterus through the cervix. It is a reversible contraceptive that can be done in a doctor's office.
Nexplanon: is about a 4 cm implant that goes into the upper forearm. This implant releases birth control hormones into the body and can last up to three years. This type of birth control has a 99% success rate for pregnancy prevention.
Dilation and curettage (D&C): an out-patient procedure to open (dilate) the cervix to collect samples of endometrial tissue with a curette. A D&C can also be done to remove a fetus that was not passed naturally after a miscarriage or to induce an abortion.
Tubal ligation: a surgery to close the fallopian tubes for the prevention of pregnancy. It is also known as "tying the tubes".
Ovarian cystectomy: the removal of a cyst that either has a solid appearance, larger than three inches in diameter, has the possibility to become cancerous, or causes a constant pain. Cysts can be removed without removing an ovary. Women who do not take birth control produce small cysts every other month but they can disappear on their own.
| Biology and health sciences | Fields of medicine | null |
52965 | https://en.wikipedia.org/wiki/Obstetrics | Obstetrics | Obstetrics is the field of study concentrated on pregnancy, childbirth and the postpartum period. As a medical specialty, obstetrics is combined with gynecology under the discipline known as obstetrics and gynecology (OB/GYN), which is a surgical field.
Main areas
Prenatal care
Prenatal care is important in screening for various complications of pregnancy. This includes routine office visits with physical exams and routine lab tests along with telehealth care for women with low-risk pregnancies:
First trimester
Routine tests in the first trimester of pregnancy generally include:
Complete blood count
Blood type
Rh-negative antenatal patients should receive RhoGAM at 28 weeks to prevent Rh disease.
Indirect Coombs test (AGT) to assess risk of hemolytic disease of the newborn
Rapid plasma reagin test to screen for syphilis
Rubella antibody screen
HBsAg test to screen for hepatitis B
Testing for chlamydia (and gonorrhea when indicated
Mantoux test for tuberculosis
Urinalysis and culture
HIV screen
Genetic screening for Down syndrome (trisomy 21) and Edwards syndrome (trisomy 18), the national standard in the United States, is rapidly evolving away from the AFP-quad screen, done typically in the second trimester at 16–18 weeks. The newer integrated screen (formerly called F.A.S.T.E.R for First And Second Trimester Early Results) can be done at 10 plus weeks to 13 plus weeks with an ultrasound of the fetal neck (thicker nuchal skin correlates with higher risk of Down syndrome being present) and two chemicals (analytes), pregnancy-associated plasma protein A and human chorionic gonadotropin (pregnancy hormone level itself). It gives an accurate risk profile very early. A second blood screen at 15 to 20 weeks refines the risk more accurately. The cost is higher than an "AFP-quad" screen due to the ultrasound and second blood test, but it is quoted to have a 93% pick up rate as opposed to 88% for the standard AFP/QS. This is an evolving standard of care in the United States.
Second trimester
MSAFP/quad. screen (four simultaneous blood tests) (maternal serum AFP, inhibin A, estriol, and βHCG) – elevations, low numbers or odd patterns correlate with neural tube defect risk and increased risks of trisomy 18 or trisomy 21
Ultrasound either abdominal or transvaginal to assess cervix, placenta, fluid and baby
Amniocentesis is the national standard for women over 35 or who reach 35 by mid pregnancy or who are at increased risk by family history or prior birth history.
Third trimester
Hematocrit (if low, the mother receives iron supplements)
Group B Streptococcus screen. If positive, the woman receives IV penicillin or ampicillin while in labor—or, if she is allergic to penicillin, an alternative therapy, such as IV clindamycin or IV vancomycin.
Glucose loading test (GLT) – screens for gestational diabetes; if > 140 mg/dL, a glucose tolerance test (GTT) is administered; a fasting glucose > 105 mg/dL suggests gestational diabetes.
Most doctors do a sugar load in a drink form of 50 grams of glucose in cola, lime or orange and draw blood an hour later (plus or minus 5 minutes). The standard modified criteria have been lowered to 135 since the late 1980s.
Fetal assessments
Obstetric ultrasonography is routinely used for dating the gestational age of a pregnancy from the size of the fetus, determine the number of fetuses and placentae, evaluate for an ectopic pregnancy and first trimester bleeding, the most accurate dating being in first trimester before the growth of the foetus has been significantly influenced by other factors. Ultrasound is also used for detecting congenital anomalies (or other foetal anomalies) and determining the biophysical profiles (BPP), which are generally easier to detect in the second trimester when the foetal structures are larger and more developed.
X-rays and computerized tomography (CT) are not used, especially in the first trimester, due to the ionizing radiation, which has teratogenic effects on the foetus. No effects of magnetic resonance imaging (MRI) on the foetus have been demonstrated, but this technique is too expensive for routine observation. Instead, obstetric ultrasonography is the imaging method of choice in the first trimester and throughout the pregnancy, because it emits no radiation, is portable, and allows for realtime imaging.
The safety of frequent ultrasound scanning has not been confirmed. Despite this, increasing numbers of women are choosing to have additional scans for no medical purpose, such as gender scans, 3D and 4D scans. A normal gestation would reveal a gestational sac, yolk sac, and fetal pole.
The gestational age can be assessed by evaluating the mean gestational sac diameter (MGD) before week 6, and the crown-rump length after week 6. Multiple gestation is evaluated by the number of placentae and amniotic sacs present.
Other tools used for assessment include:
Fetal screening is used to help assess the viability of the fetus, as well as congenital abnormalities.
Fetal karyotype can be used for the screening of genetic diseases. This can be obtained via amniocentesis or chorionic villus sampling (CVS)
Foetal haematocrit for the assessment of foetal anemia, Rh isoimmunization, or hydrops can be determined by percutaneous umbilical blood sampling (PUBS), which is done by placing a needle through the abdomen into the uterus and taking a portion of the umbilical cord.
Fetal lung maturity is associated with how much surfactant the fetus is producing. Reduced production of surfactant indicates decreased lung maturity and is a high risk factor for infant respiratory distress syndrome. Typically a lecithin:sphingomyelin ratio greater than 1.5 is associated with increased lung maturity.
Nonstress test (NST) for fetal heart rate
Oxytocin challenge test
Diseases in pregnancy
A pregnant woman may have a pre-existing disease, that may become worse or become a risk to the pregnancy, or to postnatal development of the offspring
Diabetes mellitus and pregnancy deals with the interactions of diabetes mellitus (not restricted to gestational diabetes) and pregnancy. Risks for the child include miscarriage, growth restriction, growth acceleration, foetal obesity (macrosomia), polyhydramnios and birth defects.
Lupus and pregnancy confers an increased rate of foetal death in utero and spontaneous abortion (miscarriage), as well as of neonatal lupus.
Thyroid disease in pregnancy can, if uncorrected, cause adverse effects on foetal and maternal well-being. The deleterious effects of thyroid dysfunction can also extend beyond pregnancy and delivery to affect neurointellectual development in the early life of the child. Demand for thyroid hormones is increased during pregnancy, and may cause a previously unnoticed thyroid disorder to worsen.
Hypercoagulability in pregnancy is the propensity of pregnant women to develop thrombosis (blood clots). Pregnancy itself is a factor of hypercoagulability (pregnancy-induced hypercoagulability), as a physiologically adaptive mechanism to prevent post partum bleeding. However, when combined with an additional underlying hypercoagulable states, the risk of thrombosis or embolism may become substantial.
Hyperemesis gravidarum in pregnancy occurs due to extreme, persistent nausea and vomiting during pregnancy. If untreated, can lead to dehydration, weight loss, and electrolyte imbalances. Most women develop nausea and vomiting during the first trimester. The cause of hyperemesis gravidarum is not known. However, it is believed to be caused by a rapidly rising blood level of a hormone, human chorionic gonadotropin (HCG), which is released by the placenta.
Preeclampsia is a condition that causes high blood pressure during pregnancy. If left untreated, it can be life-threatening. In pregnant women, preeclampsia may occur after 20 weeks of pregnancy, often in women who have no history of high blood pressure. Symptoms of preeclampsia may include severe headache, vision changes and pain under the ribs. However, in some women, symptoms may not occur, until they go for a routine prenatal visit.
Induction and labour
Induction is a method of artificially or prematurely stimulating labour in a woman. Reasons to induce can include pre-eclampsia, foetal distress, placental malfunction, intrauterine growth retardation and failure to progress through labour increasing the risk of infection and foetal distresses.
Induction may be achieved via several methods:
Disturbance of cervical membranes
Pessary of Prostin cream, prostaglandin E2
Intravaginal or oral administration of misoprostol
Cervical insertion of a 30-mL Foley catheter
Rupturing the amniotic membranes
Intravenous infusion of synthetic oxytocin (Pitocin or Syntocinon)
During labour, the obstetrician carries out the following tasks:
Monitor the progress of labour, by reviewing the nursing chart, performing vaginal examination, and assessing the trace produced by a foetal monitoring device (the cardiotocograph)
Provide pain relief, either by nitrous oxide, opiates, or by epidural anaesthesia done by anaesthestists, an anaesthesiologist, or a nurse anaesthetist.
Caesarean section, if there is an associated risk with vaginal delivery, as such foetal or maternal compromise.
Complications and emergencies
The main emergencies include:
Ectopic pregnancy is when an embryo implants in the uterine (fallopian) tube or (rarely) on the ovary or inside the peritoneal cavity. This may cause massive internal bleeding.
Pre-eclampsia is a disease defined by a combination of signs and symptoms that are related to maternal hypertension. The cause is unknown, and markers are being sought to predict its development from the earliest stages of pregnancy. Some unknown factors cause vascular damage in the endothelium, causing hypertension. If severe, it progresses to eclampsia, where seizures occur, which can be fatal. Preeclamptic patients with the HELLP syndrome show liver failure and disseminated intravascular coagulation (DIC). The only treatment is to deliver the foetus. Women may still develop pre-eclampsia following delivery.
Placental abruption is where the placenta detaches from the uterus and the woman and foetus can bleed to death if not managed appropriately.
Foetal distress where the foetus is getting compromised in the uterine environment.
Shoulder dystocia where one of the foetus' shoulders becomes stuck during vaginal birth. There are many risk factors, including macrosmic (large) foetus, but many are also unexplained.
Uterine rupture can occur during obstructed labour and endanger foetal and maternal life.
Prolapsed cord can only happen after the membranes have ruptured. The umbilical cord delivers before the presenting part of the foetus. If the foetus is not delivered within minutes, or the pressure taken off the cord, the foetus dies.
Obstetrical hemorrhage may be due to a number of factors such as placenta previa, uterine rupture or tears, uterine atony, retained placenta or placental fragments, or bleeding disorders.
Puerperal sepsis is an ascending infection of the genital tract. It may happen during or after labour. Signs to look out for include signs of infection (pyrexia or hypothermia, raised heart rate and respiratory rate, reduced blood pressure), and abdominal pain, offensive lochia (blood loss) increased lochia, clots, diarrhea and vomiting.
Postpartum period
The World Health Organization makes a distinction between the use of postpartum care when it concerns the care of the mother after giving birth, and postnatal care when the care of the newborn is concerned.
Postpartum care is provided to the mother following childbirth.
A woman in the Western world who gives birth in a hospital may leave the hospital as soon as she is medically stable, and chooses to leave, which can be as early as a few hours later, but usually averages a stay of one or two days; the average postnatal stay following delivery by caesarean section is three to four days.
During this time the mother is monitored for bleeding, bowel and bladder function, and baby care. The infant's health is also monitored.
Veterinary obstetrics
History
Prior to the 18th century, caring for pregnant women in Europe was confined exclusively to women, and rigorously excluded men. The expectant mother would invite close female friends and family members to her home to keep her company during childbirth. Skilled midwives managed all aspects of the labour and delivery. The presence of physicians and surgeons was very rare and only occurred if a serious complication had taken place and the midwife had exhausted all measures at her disposal. Calling a surgeon was very much a last resort and having men deliver women in this era was seen as offending female modesty.
Before the 18th century
Prior to the 18th and 19th centuries, midwifery was well established but obstetrics was not recognized as a specific medical specialty. However, the subject matter and interest in the female reproductive system and sexual practice can be traced back to Ancient Egypt and Ancient Greece. Soranus of Ephesus sometimes is called the most important figure in ancient gynecology. Living in the late first century AD and early second century, he studied anatomy and had opinions and techniques on abortion, contraception – most notably coitus interruptus – and birth complications. After his death, techniques and works of gynecology declined; very little of his works were recorded and survived to the late 18th century when gynecology and obstetrics reemerged as a medical specialism.
18th century
The 18th century marked the beginning of many advances in European midwifery, based on better knowledge of the physiology of pregnancy and labour. By the end of the century, medical professionals began to understand the anatomy of the uterus and the physiological changes that take place during labour. The introduction of forceps in childbirth also took place at this time. All these medical advances in obstetrics were a lever for the introduction of men into an arena previously managed and run by women – midwifery.
The addition of the male-midwife (or man-midwife) is historically a significant change to the profession of obstetrics. In the 18th century medical men began to train in area of childbirth and believed with their advanced knowledge in anatomy that childbirth could be improved. In France these male-midwives were referred to as accoucheurs, a title later used all over Europe. The founding of lying-in hospitals also contributed to the medicalization and male-dominance of obstetrics. These early maternity hospitals were establishments where women would come to have their babies delivered, as opposed to the practice since time immemorial of the midwife attending the home of the woman in labour. This institution provided male-midwives with endless patients to practice their techniques on and was a way for these men to demonstrate their knowledge.
Many midwives of the time bitterly opposed the involvement of men in childbirth. Some male practitioners also opposed the involvement of medical men like themselves in midwifery and even went as far as to say that male-midwives only undertook midwifery solely for perverse erotic satisfaction. The accoucheurs argued that their involvement in midwifery was to improve the process of childbirth. These men also believed that obstetrics would forge ahead and continue to strengthen.
19th century
18th-century physicians expected that obstetrics would continue to grow, but the opposite happened. Obstetrics entered a stage of stagnation in the 19th century, which lasted until about the 1880s. The central explanation for the lack of advancement during this time was the rejection of obstetrics by the medical community. The 19th century marked an era of medical reform in Europe and increased regulation over the profession. Major European institutions such as The College of Physicians and Surgeons considered delivering babies ungentlemanly work and refused to have anything to do with childbirth as a whole. Even when Medical Act 1858 was introduced, which stated that medical students could qualify as doctors, midwifery was entirely ignored. This made it nearly impossible to pursue an education in midwifery and also have the recognition of being a doctor or surgeon. Obstetrics was pushed to the side.
By the late 19th century, the foundation of modern-day obstetrics and midwifery began developing. Delivery of babies by doctors became popular and readily accepted, but midwives continued to play a role in childbirth. Midwifery also changed during this era due to increased regulation and the eventual need for midwives to become certified. Many European countries by the late 19th century were monitoring the training of midwives and issued certification based on competency. Midwives were no longer uneducated in the formal sense.
As midwifery began to develop, so did the profession of obstetrics near the end of the century. Childbirth was no longer unjustifiably despised by the medical community as it once had been at the beginning of the century. But obstetrics was underdeveloped compared to other medical specialities. Many male physicians would deliver children but very few would have referred to themselves as obstetricians. The end of the 19th century did mark a significant accomplishment in the profession with the advancements in asepsis and anaesthesia, which paved the way for the mainstream introduction and later success of the Caesarean section.
Before the 1880s mortality rates in lying-hospitals would reach unacceptably high levels and became an area of public concern. Much of these maternal deaths were due to puerperal fever, then known as childbed fever. In the 1800s Ignaz Semmelweis noticed that women giving birth at home had a much lower incidence of childbed fever than those giving birth by physicians in lying-hospitals. His investigation discovered that washing hands with an antiseptic solution before a delivery reduced childbed fever fatalities by 90%. So it was concluded that it was physicians who had been spreading disease from one labouring mother to the next. Despite the publication of this information, doctors still would not wash. It was not until the 20th century when advancements in aseptic technique and the understanding of disease would play a significant role in the decrease of maternal mortality rates among many populations.
History of obstetrics in America
The development of obstetrics as a practice for accredited doctors happened at the turn of the 18th century and thus was very differently developed in Europe and in the Americas due to the independence of many countries in the Americas from European powers. "Unlike in Europe and the British Isles, where midwifery laws were national, in America, midwifery laws were local and varied widely".
Gynaecology and Obstetrics gained attention in the American medical field at the end of the nineteenth century through the development of such procedures as the ovariotomy. These procedures then were shared with European surgeons who replicated the surgeries. This was a period when antiseptic, aseptic or anaesthetic measures were just being introduced to surgical and observational procedures and without these procedures surgeries were dangerous and often fatal. Following are two surgeons noted for their contributions to these fields include Ephraim McDowell and J. Marion Sims.
Ephraim McDowell developed a surgical practice in 1795 and performed the first ovariotomy in 1809 on a 47-year-old widow who then lived on for 31 more years. He had attempted to share this with John Bell whom he had practiced under who had retired to Italy. Bell was said to have died without seeing the document but it was published by an associate in Extractions of Diseased Ovaria in 1825. By the mid-century the surgery was both successfully and unsuccessfully being performed. Pennsylvanian surgeons the Attlee brothers made this procedure very routine for a total of 465 surgeries – John Attlee performed 64 successfully of 78 while his brother William reported 387 – between the years of 1843 and 1883. By the middle of the nineteenth century this procedure was successfully performed in Europe by English surgeons Sir Spencer Wells and Charles Clay as well as French surgeons Eugène Koeberlé, Auguste Nélaton and Jules Péan.
J. Marion Sims was the surgeon responsible for being the first treating a vesicovaginal fistula – a condition linked to many caused mainly by prolonged pressing of the foetus against the pelvis or other causes such as rape, hysterectomy, or other operations – and also having been doctor to many European royals and the 20th President of the United States James A. Garfield after he had been shot. Sims does have a controversial medical past. Under the beliefs at the time about pain and the prejudice towards African people, he had practiced his surgical skills and developed skills on slaves. These women were the first patients of modern gynecology. One of the women he operated on was named Anarcha Westcott, the woman he first treated for a fistula.
Historical role of gender
Women and men inhabited very different roles in natal care up to the 18th century. The role of a physician was exclusively held by men who went to university, an overly male institution, who would theorize anatomy and the process of reproduction based on theological teaching and philosophy. Many beliefs about the female body and menstruation in the 17th and 18th centuries were inaccurate; clearly resulting from the lack of literature about the practice. Many of the theories of what caused menstruation prevailed from Hippocratic philosophy. Midwives, meaning "with woman", were those who assisted in the birth and care of both born and unborn children, a position historically held mainly by women.
During the birth of a child, men were rarely present. Women from the neighbourhood or family would join in on the process of birth and assist in many different ways. The one position where men would help with the birth of a child would be in the sitting position, usually when performed on the side of a bed to support the mother.
Men entered the field of obstetrics in the nineteenth century, resulting in a change of focus within the profession. Gynecology developed as a new and separate field of study from obstetrics, focusing on the curing of illness and indispositions of female sexual organs, encompassing conditions such as menopause, uterine and cervical problems, and tissue damage as a result of childbirth.
| Biology and health sciences | Fields of medicine | Health |
52967 | https://en.wikipedia.org/wiki/Gynaecology | Gynaecology | Gynaecology or gynecology (see American and British English spelling differences) is the area of medicine that involves the treatment of women's diseases, especially those of the female reproductive organs. It is often paired with the field of obstetrics, which focuses on pregnancy and childbirth, thereby forming the combined area of obstetrics and gynaecology (OB-GYN).
The term comes from Greek and means . Its counterpart is andrology, which deals with medical issues specific to the male reproductive system.
Etymology
The word gynaecology comes from the oblique stem () of the Greek word () meaning , and meaning .
History
Antiquity
The Kahun Gynaecological Papyrus, dated to about 1800 BC, deals with gynaecological diseases, fertility, pregnancy, contraception, etc. The text is divided into thirty-four sections, each section dealing with a specific problem and containing diagnosis and treatment; no prognosis is suggested. Treatments are non-surgical, comprising applying medicines to the affected body part or swallowing them. The womb is at times seen as the source of complaints manifesting themselves in other body parts.
Ayurveda, an Indian traditional medical system, also provides details about concepts and techniques related to gynaecology.
The Hippocratic Corpus contains several gynaecological treatises dating to the 5th and 4th centuries BC. Aristotle is another strong source for medical texts from the 4th century BC with his descriptions of biology primarily found in History of Animals, Parts of Animals, Generation of Animals. The gynaecological treatise Gynaikeia by Soranus of Ephesus (1st/2nd century AD) is extant (together with a 6th-century Latin paraphrase by Muscio, a physician of the same school). He was the chief representative of the school of physicians known as the "methodists".
Modern gynaecology. J. Marion Sims
In the medical schools of the early nineteenth century, doctors did not study female reproductive anatomy, seen as repulsive, nor train in pregnancy and childbirth management. That women, because of their anatomy and the risks of the dangerous birthing process, had unique medical concerns and challenges, enough that a doctor might specialize in them, is an innovation widely credited to J. Marion Sims and to a lesser extent his trainee and partner Nathan Bozeman, physicians from Montgomery, Alabama. Sims is widely considered to be the father of modern gynaecology. While there have been isolated precedents for some of his innovations, he was the first to have published on the Sims' position, the Sims' speculum, the Sims sigmoid catheter, and on gynecological surgery, first on repair of vesico-vaginal fistulas, a socially devastating consequence of protracted childbirth, at the time without treatment of any sort. He founded the first women's hospital in the country, first in his backyard in Montgomery, limited to Black enslaved women, then the Woman's Hospital of New York.
He was elected president of the American Medical Association, and was the first American physician of whom a statue was erected.
Sims developed his new specialty using the bodies of enslaved women, who could not refuse the extended glance of any white male that cared to observe any part of their anatomy. They could not "consent" in the sense modern medical research requires.
At the time anesthesia was itself a research area, and the first experiments (in dentistry) were being published. Using early anesthesia (in 1845, say) was much more dangerous and difficult than it would be a century later. In addition, it was widely believed that Blacks did not feel pain as much as whites, and white women proved unable to endure the pain.
At the time, Sims was seen as a hero. Even his enemies, Bozeman chief among them, did not attack him for either experimenting on the enslaved, or for not using anesthesia. Abolitionists such as William Lloyd Garrison were quick to put in print any mistreatment of the enslaved; Garrison's influential The Liberator has been completely indexed, but it never mentions Sims. Nor does the digitized portion of the Black press mention him. When he left Alabama in 1853, a local newspaper called him "an honor to our state".
In the late 20th century, Sims has come to be villainized. Now criticized for his practices, Sims developed some of his techniques and instruments by operating on slaves, many of whom were not given anesthesia. Sims performed surgeries on 12 enslaved women in his homemade backyard hospital for four years. While performing these surgeries he invited eager physicians and students to watch invasive and painful procedures while the women were exposed. On one of the women, named Anarcha, he performed 30 surgeries without anesthesia. Due to having so many enslaved women, he would rotate from one to another, continuously trying to perfect the repair of their fistulas. Physicians and students lost interest in assisting Sims over the course of his backyard practice, and he recruited other enslaved women, who were healing from their own surgeries, to assist him. In 1855, Sims went on to found the Woman's Hospital in New York, the first hospital specifically for female disorders.
Examination
In some countries, women must first see a general practitioner (GP; also known as a family practitioner (FP)) prior to seeing a gynaecologist. If their condition requires training, knowledge, surgical procedure, or equipment unavailable to the GP, the patient is then referred to a gynaecologist. In other countries, laws may allow patients to see gynaecologists without a referral. Some gynaecologists provide primary care in addition to aspects of their own specialty. With this option available, some women opt to see a gynaecological surgeon for non-gynaecological problems without another physician's referral.
As in all of medicine, the main tools of diagnosis are clinical history, examination and investigations. Gynaecological examination is quite intimate, more so than a routine physical exam. It also requires unique instrumentation such as the speculum. The speculum consists of two hinged blades of concave metal or plastic which are used to retract the tissues of the vagina and permit examination of the cervix, the lower part of the uterus located within the upper portion of the vagina. Gynaecologists typically do a bimanual examination (one hand on the abdomen and one or two fingers in the vagina) to palpate the cervix, uterus, ovaries and bony pelvis. It is not uncommon to do a rectovaginal examination for a complete evaluation of the pelvis, particularly if any suspicious masses are appreciated. Male gynaecologists may have a female chaperone for their examination. An abdominal or vaginal ultrasound can be used to confirm any abnormalities appreciated with the bimanual examination or when indicated by the patient's history.
Diseases
Examples of conditions dealt with by a gynaecologist are:
Cancer and pre-cancerous diseases of the reproductive organs including ovaries, fallopian tubes, uterus, cervix, vagina, and vulva
Incontinence of urine
Amenorrhoea (absent menstrual periods)
Endometriosis
Dysmenorrhoea (painful menstrual periods)
Infertility
Menorrhagia (heavy menstrual periods); a common indication for hysterectomy
Prolapse of pelvic organs
Infections of the vagina (vaginitis), cervix and uterus (including fungal, bacterial, viral, and protozoal)
Pelvic inflammatory disease
Urinary tract infections
Polycystic ovary syndrome
Premenstrual dysphoric disorder
Post-menopausal osteoporosis
Other vaginal diseases
There is some crossover in these areas. For example, a woman with urinary incontinence may be referred to a urologist.
Therapies
As with all surgical specialties, gynaecologists may employ medical or surgical therapies (or many times, both), depending on the exact nature of the problem that they are treating. Pre- and post-operative medical management will often employ many standard drug therapies, such as antibiotics, diuretics, antihypertensives, and antiemetics. Additionally, gynaecologists make frequent use of specialized hormone-modulating therapies (such as Clomifene citrate and hormonal contraception) to treat disorders of the female genital tract that are responsive to pituitary or gonadal signals.
Surgery, however, is the mainstay of gynaecological therapy. For historical and political reasons, gynaecologists were previously not considered "surgeons", although this point has always been the source of some controversy. Modern advancements in both general surgery and gynaecology, however, have blurred many of the once rigid lines of distinction. The rise of sub-specialties within gynaecology which are primarily surgical in nature (for example urogynaecology and gynaecological oncology) have strengthened the reputations of gynaecologists as surgical practitioners, and many surgeons and surgical societies have come to view gynaecologists as comrades of sorts. As proof of this changing attitude, gynaecologists are now eligible for fellowship in both the American College of Surgeons and Royal Colleges of Surgeons, and many newer surgical textbooks include chapters on (at least basic) gynaecological surgery.
Some of the more common operations that gynaecologists perform include:
Dilation and curettage (removal of the uterine contents for various reasons, including completing a partial miscarriage and diagnostic sampling for dysfunctional uterine bleeding refractive to medical therapy)
Hysterectomy (removal of the uterus)
Oophorectomy (removal of the ovaries)
Tubal ligation (a type of permanent sterilization)
Hysteroscopy (inspection of the uterine cavity)
Diagnostic laparoscopy – used to diagnose and treat sources of pelvic and abdominal pain. Laparoscopy is the only way to accurately diagnose pelvic/abdominal endometriosis.
Exploratory laparotomy – may be used to investigate the level of progression of benign or malignant disease, or to assess and repair damage to the pelvic organs.
Various surgical treatments for urinary incontinence, including cystoscopy and sub-urethral slings.
Surgical treatment of pelvic organ prolapse, including correction of cystocele and rectocele.
Appendectomy – often performed to remove site of painful endometriosis implantation or prophylactically (against future acute appendicitis) at the time of hysterectomy or Caesarean section. May also be performed as part of a staging operation for ovarian cancer.
Cervical Excision Procedures (including cryosurgery) – removal of the surface of the cervix containing pre-cancerous cells which have been previously identified on Pap smear.
Specialist training
In the UK the Royal College of Obstetricians and Gynaecologists, based in London, encourages the study and advancement of both the science and practice of obstetrics and gynaecology. This is done through postgraduate medical education and training development, and the publication of clinical guidelines and reports on aspects of the specialty and service provision. The RCOG International Office works with other international organisations to help lower maternal morbidity and mortality in under-resourced countries.
Gynaecologic oncology is a subspecialty of gynaecology, dealing with gynaecology-related cancer.
Urogynaecology is a subspecialty of gynaecology and urology dealing with urinary or fecal incontinence and pelvic organ prolapse.
Gender of physicians
Improved access to education and the professions in recent decades has seen women gynaecologists outnumber men in the once male-dominated medical field of gynaecology. In some gynaecological sub-specialties, where an over-representation of males persists, income discrepancies appear to show male practitioners earning higher averages.
Speculations on the decreased numbers of male gynaecologist practitioners report a perceived lack of respect from within the medical profession, limited future employment opportunities and questions to the motivations and character of men who choose the medical field concerned with female sexual organs.
Surveys of women's views on the issue of male doctors conducting intimate examinations show a large and consistent majority found it uncomfortable, were more likely to be embarrassed and less likely to talk openly or in detail about personal information, or discuss their sexual history with a man. The findings raised questions about the ability of male gynaecologists to offer quality care to patients. This, when coupled with more women choosing female physicians has decreased the employment opportunities for men choosing to become gynaecologists.
In the United States, it has been reported that four in five students choosing a residency in gynaecology are now female. In several places in Sweden, to comply with discrimination laws, patients may not choose a doctor—regardless of specialty—based on factors such as ethnicity or gender and declining to see a doctor solely because of preference regarding e.g. the practitioner's skin color or gender may legally be viewed as refusing care. In Turkey, due to patient preference to be seen by another female, there are now few male gynaecologists working in the field.
There have been a number of legal challenges in the US against healthcare providers who have started hiring based on the gender of physicians. Mircea Veleanu argued, in part, that his former employers discriminated against him by accommodating the wishes of female patients who had requested female doctors for intimate exams. A male nurse complained about an advert for an all-female obstetrics and gynaecology practice in Columbia, Maryland, claiming this was a form of sexual discrimination. In 2000, David Garfinkel, a New Jersey-based OB-GYN, sued his former employer after being fired due to, as he claimed, "because I was male, I wasn't drawing as many patients as they'd expected".
| Biology and health sciences | Fields of medicine | Health |
52974 | https://en.wikipedia.org/wiki/Emergency%20medicine | Emergency medicine | Emergency medicine is the medical specialty concerned with the care of illnesses or injuries requiring immediate medical attention. Emergency medicine physicians (often called "ER doctors" in the United States) specialize in providing care for unscheduled and undifferentiated patients of all ages. As first-line providers, in coordination with emergency medical services, they are primarily responsible for initiating resuscitation and stabilization and performing the initial investigations and interventions necessary to diagnose and treat illnesses or injuries in the acute phase. Emergency medical physicians generally practice in hospital emergency departments, pre-hospital settings via emergency medical services, and intensive care units. Still, they may also work in primary care settings such as urgent care clinics.
Sub-specializations of emergency medicine include; disaster medicine, medical toxicology, point-of-care ultrasonography, critical care medicine, emergency medical services, hyperbaric medicine, sports medicine, palliative care, or aerospace medicine.
Various models for emergency medicine exist internationally. In countries following the Anglo-American model, emergency medicine initially consisted of surgeons, general practitioners, and other generalist physicians. However, in recent decades it has become recognised as a specialty in its own right with its training programmes and academic posts, and the specialty is now a popular choice among medical students and newly qualified medical practitioners. By contrast, in countries following the Franco-German model, the specialty does not exist, and emergency medical care is instead provided directly by anesthesiologists (for critical resuscitation), surgeons, specialists in internal medicine, paediatricians, cardiologists or neurologists as appropriate. Emergency medicine is still evolving in developing countries, and international emergency medicine programs offer hope of improving primary emergency care where resources are limited.
Scope
Emergency medicine is a medical specialty—a field of practice based on the knowledge and skills required to prevent, diagnose, and manage acute and urgent aspects of illness and injury affecting patients of all age groups with a full spectrum of undifferentiated physical and behavioural disorders. It further encompasses an understanding of the development of pre-hospital and in-hospital emergency medical systems and the skills necessary for this development.
The field of emergency medicine encompasses care involving the acute care of internal medical and surgical conditions. In many modern emergency departments, emergency physicians see many patients, treating their illnesses and arranging for disposition—either admitting them to the hospital or releasing them after treatment as necessary. They also provide episodic primary care to patients during off-hours and those who do not have primary care providers. Most patients present to emergency departments with low-acuity conditions (such as minor injuries or exacerbations of chronic disease), but a small proportion will be critically ill or injured. Therefore, the emergency physician requires broad knowledge and procedural skills, often including surgical procedures, trauma resuscitation, advanced cardiac life support and advanced airway management. They must have some of the core skills from many medical specialities—the ability to resuscitate a patient (intensive care medicine), manage a difficult airway (anesthesiology), suture a complex laceration (plastic surgery), set a fractured bone or dislocated joint (orthopaedic surgery), treat a heart attack (cardiology), manage strokes (neurology), work-up a pregnant patient with vaginal bleeding (obstetrics and gynaecology), control a patient with mania (psychiatry), stop a severe nosebleed (otolaryngology), place a chest tube (cardiothoracic surgery), and conduct and interpret x-rays and ultrasounds (radiology). This generalist approach can obviate barrier-to-care issues seen in systems without specialists in emergency medicine, where patients requiring immediate attention are instead managed from the outset by specialty doctors such as surgeons or internal physicians. However, this may lead to barriers through acute and critical care specialities disconnecting from emergency care.
Emergency medicine may separate from urgent care, which refers to primary healthcare for less emergent medical issues, but there is obvious overlap, and many emergency physicians work in urgent care settings. Emergency medicine also includes many aspects of acute primary care and shares with family medicine the uniqueness of seeing all patients regardless of age, gender or organ system. The emergency physician workforce also includes many competent physicians who have medical skills from other specialities.
Physicians specializing in emergency medicine can enter fellowships to receive credentials in subspecialties such as palliative care, critical care medicine, medical toxicology, wilderness medicine, pediatric emergency medicine, sports medicine, disaster medicine, tactical medicine, ultrasound, pain medicine, pre-hospital emergency medicine, or undersea and hyperbaric medicine.
The practice of emergency medicine is often quite different in rural areas where there are far fewer other specialities and healthcare resources. In these areas, family physicians with additional skills in emergency medicine often staff emergency departments. Rural emergency physicians may be the only health care providers in the community and require skills that include primary care and obstetrics.
Work patterns
Patterns vary by country and region. In the United States, the employment arrangement of emergency physician practices are either private (with a co-operative group of doctors staffing an emergency department under contract), institutional (physicians with or without an independent contractor relationship with the hospital), corporate (physicians with an independent contractor relationship with a third-party staffing company that services multiple emergency departments), or governmental (for example, when working within personal service military services, public health services, veterans' benefit systems or other government agencies).
In the United Kingdom, all consultants in emergency medicine work in the National Health Service, and there is little scope for private emergency practice. In other countries like Australia, New Zealand, or Turkey, emergency medicine specialists are almost always salaried employees of government health departments and work in public hospitals, with pockets of employment in private or non-government aeromedical rescue or transport services, as well as some private hospitals with emergency departments; they may be supplemented or backed by non-specialist medical officers, and visiting general practitioners. Rural emergency departments are sometimes run by general practitioners alone, sometimes with non-specialist qualifications in emergency medicine.
History
During the French Revolution, after seeing the speed with which the carriages of the French flying artillery maneuvered across the battlefields, French military surgeon Dominique Jean Larrey applied the idea of ambulances, or "flying carriages", for rapid transport of wounded soldiers to a central place where medical care was more accessible and practical. Larrey operated ambulances with trained crews of drivers, corpsmen and litter-bearers and had them bring the wounded to centralized field hospitals, effectively creating a forerunner of the modern MASH units. Dominique Jean Larrey is sometimes called the Father of Emergency Medicine for his strategies during the French wars.
Emergency medicine as an independent medical specialty is relatively young. Before the 1960s and 1970s, hospital emergency departments (EDs) were generally staffed by physicians on staff at the hospital on a rotating basis, among them family physicians, general surgeons, internists, and a variety of other specialists. In many smaller emergency departments, nurses would triage patients, and physicians would be called in based on the type of injury or illness. Family physicians were often on call for the emergency department and recognized the need for dedicated emergency department coverage. Many of the pioneers of emergency medicine were family physicians and other specialists who saw a need for additional training in emergency care.
During this period, physicians began to emerge who had left their respective practices to devote their work entirely to the ED. In the UK in 1952, Maurice Ellis was appointed as the first "casualty consultant" at Leeds General Infirmary. In 1967, the Casualty Surgeons Association was co-established with Maurice Ellis as its first president. In the US, the first of such groups managed by Dr James DeWitt Mills in 1961, along with four associate physicians; Dr Chalmers A. Loughridge, Dr William Weaver, Dr John McDade, and Dr Steven Bednar, at Alexandria Hospital in Alexandria, Virginia, established 24/7 year-round emergency care, which became known as the "Alexandria Plan".
It was not until Dr. John Wiegenstein founded the American College of Emergency Physicians (ACEP) the recognition of emergency medicine training programs by the AMA and the AOA, and in 1979 a historic vote by the American Board of Medical Specialties that emergency medicine became a recognized medical specialty in the US. The first emergency medicine residency program in the world began in 1970 at the University of Cincinnati. Furthermore, the first department of emergency medicine at a US medical school occurred in 1971 at the University of Southern California. The second residency program in the United States soon followed at what was then called Hennepin County General Hospital in Minneapolis, with two residents entering the program in 1971.
In 1990 the UK's Casualty Surgeons Association changed its name to the British Association for Accident and Emergency Medicine and subsequently became the British Association for Emergency Medicine (BAEM) in 2004. In 1993, an intercollegiate Faculty of Accident and Emergency Medicine (FAEM) became a "daughter college" of six royal medical colleges in England and Scotland to arrange professional examinations and training. In 2005, the BAEM and the FAEM became a single unit to form the College of Emergency Medicine, now the Royal College of Emergency Medicine, which conducts membership and fellowship examinations and publishes guidelines and standards for the practice of emergency medicine.
Financing and practice organization
Reimbursement
Many hospitals and care centres feature departments of emergency medicine, where patients can receive acute care without an appointment. While many patients get treated for life-threatening injuries, others utilize the emergency department (ED) for non-urgent reasons such as headaches or a cold. (defined as "visits for conditions for which a delay of several hours would not increase the likelihood of an adverse outcome"). As such, EDs can adjust staffing ratios and designate an area of the department for faster patient turnover to accommodate various patient needs and volumes. Policies have improved to assist better ED staff (such as emergency medical technicians, paramedics). The emergency department, welfare programs, and healthcare clinics serve as a critical part of the healthcare safety net for uninsured patients who cannot afford medical treatment or adequately utilize their coverage.
In emergency departments in Australia, the government utilises an "Activity based funding and management", meaning that the amount of funding to emergency departments are allocated money based on the number of patients and the complexity of their cases or illnesses. However, rural emergency departments of Australia are funded under the principle of providing the necessary equipment and staffing levels required to provide safe and adequate care, not necessarily on the number of patients.
Compensation
In the United States, Emergency Physicians are compensated at a higher rate than some other specialities, ranking 10th out of 26 physician specialities in 2015, at an average salary of $306,000 annually. They are compensated in the mid-range (averaging $13,000 annually) for non-patient activities, such as speaking engagements or acting as an expert witness; they also saw a 12% increase in salary from 2014 – 2015 (which was not out of line with many other physician specialities that year). While emergency physicians work 8–12 hour shifts and do not tend to work on-call, the high level of stress and need for solid diagnostic and triage capabilities for the undifferentiated, acute patient contributes to arguments justifying higher salaries for these physicians. Emergency care must be available every hour of every day and requires a doctor to be available on-site 24/7, unlike an outpatient clinic or other hospital departments with more limited hours and may only call a physician in when needed. The necessity to have a physician on staff and all other diagnostic services available every hour of every day is thus a costly arrangement for hospitals.
Payment systems
American health payment systems are undergoing significant reform efforts, Which include compensating emergency physicians through "pay for performance" incentives and penalty measures under commercial and public health programs, including Medicare and Medicaid. This payment reform aims to improve the quality of care and control costs, despite the differing opinions on the existing evidence to show that this payment approach is effective in emergency medicine. Initially, these incentives would only target primary care providers (PCPs), but some would argue that emergency medicine is primary care, as no one refers patients to the ED. In one such program, two specific conditions listed were directly tied to patients frequently seen by emergency medical providers: acute myocardial infarction and pneumonia.(See: Hospital Quality Incentive Demonstration.)
There are some challenges with implementing these quality-based incentives in emergency medicine in that patients are often not given a definitive diagnosis in the ED, making it challenging to allocate payments through coding. Additionally, adjustments based on patient risk-level and multiple co-morbidities for complex patients further complicate attribution of positive or negative health outcomes. It is not easy to assess whether much of the costs directly result from the emergent condition treated in acutely care settings. It is also difficult to quantify the savings due to preventive care during emergency treatment (i.e. workup, stabilizing treatments, coordination of care and discharge, rather than a hospital admission). Thus, ED providers tend to support a modified fee-for-service model over other payment systems.
Overutilization
Some patients without health insurance utilize EDs as their primary form of medical care, as their financial status limits their access to consistent care. Because these patients cannot utilize insurance or primary care systems, emergency medical providers often increased volumes of lower acuity patients and risk of financial loss, especially since many patients cannot pay for their care (see below). ED overuse produces $38 billion in spending each year (i.e. care delivery and coordination failures, over-treatment, administrative complexity, pricing failures, and fraud), Moreover, it unnecessarily drains departmental resources, reducing the quality of care across all patients. While overuse is not limited to the uninsured, the uninsured constitute a growing proportion of non-urgent ED visits. Insurance coverage can help mitigate overutilization by improving access to alternative forms of care and lowering the need for emergency visits.
A common misconception identifies frequent ED visitors as a significant factor in excess spending. However, frequent ED users make up a small portion of those contributing to overutilization and are often insured.
Uncompensated care
Injury and illness are often unforeseen, and patients of lower socioeconomic status are especially susceptible to being suddenly burdened with the cost of a necessary ED visit. For example, in the event that a patient is unable to pay for medical care received, the hospital, under the Emergency Medical Treatment and Active Labor Act (EMTALA), is obligated to treat emergency conditions regardless of a patient's ability to pay and therefore faces an economic loss for this uncompensated care. Estimates suggest that over half (approximately 55%) of all quantifiable emergency care is uncompensated and inadequate reimbursement has led to the closure of many EDs. Policy changes (such as the Affordable Care Act) are expected to decrease the number of uninsured people and thereby reduce uncompensated care.
In addition to decreasing the uninsured rate, ED overutilization might reduce by improving patient access to primary care and increasing patient flow to alternative care centres for non-life-threatening injuries. Financial disincentives, patient education, and improved management for patients with chronic diseases can also reduce overutilization and help manage costs of care. Moreover, physician knowledge of prices for treatment and analyses, discussions on costs with their patients, and a changing culture away from defensive medicine can improve cost-effective use. A transition towards more value-based care in the ED is an avenue by which providers can contain costs.
EMTALA
Doctors that work in the EDs of hospitals receiving Medicare funding are subject to the provisions of EMTALA. The US Congress enacted EMTALA in 1986 to curtail "patient dumping", a practice whereby patients were refused medical care for economic or other non-medical reasons. Since its enactment, ED visits have substantially increased, with one study showing a rise in visits of 26% (which is more than double the increase in population over the same period). While more individuals are receiving care, a lack of funding and ED overcrowding may be affecting quality. To comply with the provisions of EMTALA, hospitals, through their ED physicians, must provide medical screening and stabilize the emergency medical conditions of anyone that presents themselves at a hospital ED with patient capacity. EMTALA holds both the hospital and the responsible ED physician liable for civil penalties of up to $50,000 if there is no help for those in need. . While both the Office of Inspector General, U.S. Department of Health and Human Services (OIG) and private citizens can bring an action under EMTALA, courts have uniformly held that ED physicians can only be held liable if the case is prosecuted by OIG (whereas hospitals are subject to penalties regardless of who brings the suit). Additionally, the Centres for Medicare and Medicaid Services (CMS) can discontinue provider status under Medicare for physicians that do not comply with EMTALA. Liability also extends to on-call physicians that fail to respond to an ED request to come to the hospital to provide service. While the goals of EMTALA are laudable, commentators have noted that it appears to have created a substantial unfunded burden on the resources of hospitals and emergency physicians. As a result of financial difficulty, between the period of 1991–2011, 12.6% of EDs in the US closed.
Care delivery in different ED settings
Rural
Despite the practice emerging over the past few decades, the delivery of emergency medicine has significantly increased and evolved across diverse settings related to cost, provider availability and overall usage. Before the Affordable Care Act (ACA), low-acuity emergency medicine visits were leveraged primarily by "uninsured or underinsured patients, women, children, and minorities, all of whom frequently face barriers to accessing primary care". While this still exists today, as mentioned above, it is critical to consider the location in which care is delivered to understand the population and system challenges related to overutilization and high cost. In rural communities where provider and ambulatory facility shortages exist, a primary care physician (PCP) in the ED with general knowledge is likely to be the only source of health care for a population, as specialists and other health resources are generally unavailable due to lack of funding and desire to serve in these areas. As a result, the incidence of complex co-morbidities not managed by the appropriate provider results in worse health outcomes and eventually costlier care that extends beyond rural communities. Though typically quite separated, PCPs in rural areas must partner with larger health systems to comprehensively address the complex needs of their community, improve population health, and implement strategies such as telemedicine to improve health outcomes and reduce ED utilization for preventable illnesses. (See: Rural health.)
Rural care has benefitted in the post-pandemic (2020) era by the rapid expansion of telemedicine programs, including those that assist with Emergency Medical care. This ha enhanced the ability of non-Emergency Medicine boarded physicians, physician assistants and nurse practitioners to provide a higher level of care by partnering with Emergency Physicians at larger centers, via telehealth.
Urban
Alternatively, emergency medicine in urban areas consists of diverse provider groups, including physicians, physician assistants, nurse practitioners and registered nurses who coordinate with specialists in both inpatient and outpatient facilities to address patients' needs, more specifically in the ED. For all systems, regardless of funding source, EMTALA mandates EDs to conduct a medical examination for anyone that presents at the department, irrespective of paying ability. Non-profit hospitals and health systems – as required by the ACA – must provide a certain threshold of charity care "by actively ensuring that those who qualify for financial assistance get it, by charging reasonable rates to uninsured patients and by avoiding extraordinary collection practices." While there are limitations, this mandate provides support to many in need. That said, despite policy efforts and increased funding and federal reimbursement in urban areas, the triple aim (of improving patient experience, enhancing population health, and reducing the per-capita cost of care) remains a challenge without providers' and payers' collaboration to increase access to preventive care and decrease in ED usage. As a result, many experts support the notion that emergency medical services should only serve immediate risks in urban and rural areas.
Patient–provider relationships
As stated above, EMTALA includes provisions that protect patients from being turned away or transferred before adequate stabilisation. Upon making contact with a patient, EMS providers are responsible for diagnosing and stabilising a patient's condition without regard for the ability to pay. In the pre-hospital setting, providers must exercise appropriate judgement in choosing a suitable hospital for transport. Hospitals can only turn away incoming ambulances if they are on diversion and incapable of providing adequate care. However, once a patient has arrived on hospital property, care must be provided. At the hospital, a triage nurse first contacts the patient, who determines the appropriate level of care needed.
According to Mead v. Legacy Health System, a patient-physician relationship is established when "the physician takes an affirmative action with regard to the care of the patient". Initiating such a relationship forms a legal contract in which the physician must continue to provide treatment or adequately terminate the relationship. This legal responsibility can extend to physician consultations and on-call physicians even without direct patient contact. In emergency medicine, termination of the patient–provider relationship prior to stabilization or without handoff to another qualified provider is considered abandonment. In order to initiate an outside transfer, a physician must verify that the next hospital can provide a similar or higher level of care. Hospitals and physicians must also ensure that the patient's condition will not be further aggravated by the transfer process.
The setting of emergency medicine presents a challenge for delivering high quality, patient-centered care. Clear, effective communication can be particularly difficult due to noise, frequent interruptions, and high patient turnover. The Society for Academic Emergency Medicine has identified five essential tasks for patient-physician communication: establishing rapport, gathering information, giving information, providing comfort, and collaboration. The miscommunication of patient information is a crucial source of medical error; minimising shortcoming in communication remains a topic of current and future research.
Medical error
Many circumstances, including the regular transfer of patients in emergency treatment and crowded, noisy and chaotic ED environments, make emergency medicine particularly susceptible to medical error and near misses. One study identified an error rate of 18 per 100 registered patients in one particular academic ED. Another study found that where a lack of teamwork (i.e. poor communication, lack of team structure, lack of cross-monitoring) was implicated in a particular incident of ED medical error, "an average of 8.8 teamwork failures occurred per case [and] more than half of the deaths and permanent disabilities that occurred were judged avoidable." Particular cultural (i.e. "a focus on the errors of others and a 'blame-and-shame' culture") and structural (i.e. lack of standardisation and equipment incompatibilities) aspects of emergency medicine often result in a lack of disclosure of medical error and near misses to patients and other caregivers. While concerns about malpractice liability are one reason why disclosure of medical errors is not made, some have noted that disclosing the error and providing an apology can mitigate malpractice risk. Ethicists uniformly agree that the disclosure of a medical error that causes harm is a care provider's duty. The critical components of the disclosure include "honesty, explanation, empathy, apology, and the chance to lessen the chance of future errors" (represented by the mnemonic HEEAL). The nature of emergency medicine is such that error will likely always be a substantial risk of emergency care. However, maintaining public trust through open communication regarding a harmful error can help patients and physicians constructively address problems when they occur.
Treatments
Emergency medicine is a primary or first-contact point of care for patients requiring the use of the health care system. Specialists in emergency medicine are required to possess specialist skills in acute illness diagnosis and resuscitation. Emergency physicians are responsible for providing immediate recognition, evaluation, care, and stabilisation to adult and pediatric patients in response to acute illness and injury.
Emergency medical physicians provide treatments to a range of cases requiring vast knowledge. They deal with patients from mental illnesses to physical and anything in-between. An average treatment process would likely involve, investigation then diagnosis then either treatment or the patient being admitted. In terms of procedure's they cover a wide and broad range, including treatment to GSW's (Gun Shot Wounds), Head and body traumas, stomach bugs, mental episodes, seizures and much more. They are some of the most highly trained physicians in the world and are responsible for providing immediate recognition, evaluation, care, and stabilisation to adult and paediatric patients in response to acute illness and injury.As well as being the first point of care for many patients in emergency situations.
Training
There are a variety of international models for emergency medicine training. There are two different models among those with well-developed training programs: a "specialist" model or "a multidisciplinary model". Additionally, in some countries, the emergency medicine specialist rides in the ambulance. For example, in France and Germany, the physician, often an anesthesiologist, rides in the ambulance and provides stabilising care at the scene. The patient is directed to the appropriate hospital department, so emergency care is much more multidisciplinary than the Anglo-American model.
In countries such as the US, the United Kingdom, Canada and Australia, ambulances crewed by paramedics and emergency medical technicians respond to out-of-hospital emergencies and transport patients to emergency departments, meaning there is more dependence on these healthcare providers and there is more dependence on paramedics and EMTs for on-scene care. Emergency physicians are therefore more "specialists" since all patients are taken to the emergency department. Most developing countries follow the Anglo-American model: the gold standard is three or four-year independent residency training programs in emergency medicine. Some countries develop training programs based on a primary care foundation with additional emergency medicine training. In developing countries, there is an awareness that Western models may not be applicable and may not be the best use of limited health care resources. For example, specialty training and pre-hospital care in developed countries are too expensive and impractical for use in many developing countries with limited health care resources. International emergency medicine provides a critical global perspective and hope for improvement in these areas.
A brief review of some of these programs follows:
Argentina
In Argentina, the SAE (Sociedad Argentina de Emergencias) is the leading organisation of emergency medicine. There are many residency programs. Also, it is possible to reach the certification with a two-year postgraduate university course after a few years of ED background.
Australia and New Zealand
The specialist medical college responsible for emergency medicine in Australia and New Zealand is the Australasian College for Emergency Medicine (ACEM). The training program is nominally seven years in duration, after which the trainee is awarded a Fellowship of ACEM, conditional upon passing all necessary assessments.
Dual fellowship programs also exist for paediatric medicine (in conjunction with the Royal Australasian College of Physicians) and intensive care medicine (in conjunction with the College of Intensive Care Medicine). These programs nominally add one or more years to the ACEM training program.
For medical doctors not (and not wishing to be) specialists in emergency medicine but have a significant interest or workload in emergency departments, the ACEM provides non-specialist certificates and diplomas.
The Australian College of Rural and Remote Medicine (ACRRM) is the responsible body for the training and upholding of standards for practice and provision of rural and remote medical care. Prospective rural generalists undertaking this four-year fellowship program have an opportunity to complete Advanced Specialised Training (AST) in emergency medicine.
Belgium
In Belgium there are three recognised ways to practice emergency medicine. Until 2005 there was no accredited emergency medicine program. Emergency medicine was performed by general practitioners (having followed a 240-hour course, Acute Medicine) or by specialists (surgeon, internal medicine, neurologist, anesthesiologist) with or without supra-specialty training in emergency medicine.
Since 2005 residency training exists for acute medicine (3 years) or emergency medicine (6 years). At least 50% of the training is in the emergency department; the other part is a rotation between disciplines like pediatrics, surgery, orthopedic surgery, anesthesiology and critical care medicine.
Alternative an attending physician with one of following specialties (anesthesiology, internal medicine, cardiology, gastro-enterology, pneumology, rheumatology, urology, general surgery, plastic & reconstructive surgery, orthopedic surgery, neurology, neurosurgery, pediatrics) can follow a supra-specialty program of two years to become an emergency medicine specialist.
Brazil
In Brazil, the first emergency medicine residency program was created at Hospital Pronto Socorro de Porto Alegre in 1996. In 2002, the emergency medical services were standardized nationally with the creation of SAMU (Serviço de atendimento móvel de urgência), inspired by French EMS, which also provides training to its employees. The nacional emergency medicina association (ABRAMEDE – Associação Brasileira de Medicina de Emergência) was created in 2007. In 2008 the second residency program was started at Messejana Hospital in Fortaleza. Then, in 2015, emergency medicine was formally recognized as a medical specialty by the Brazilian Medical Association. After formal recognition, multiple residency programs were created nationwide (e.g. Universidade Federal de Minas Gerais in 2016 and Universidade de São Paulo in 2017). The residency consists of a three-year program with training in all emergency department specialties (i.e. internal medicine, surgery, pediatrics, orthopedics, OB/GYN), EMS and intensive care.
Chile
In Chile, emergency medicine begins its journey in Chile with the first specialty program at the beginning of the 1990s, at the University of Chile and the University of Santiago of Chile. Currently, it is a primary specialty legally recognised by the Ministry of Health since 2013. It has multiple training programs for specialists, notably the University of Chile, Pontifical Catholic University of Chile, Clínica Alemana – Universidad del Desarrollo, San Sebastian University – MUE and University of Santiago of Chile (USACH). Currently, and intending to strengthen the specialty at the country level, FOAMed initiatives have emerged (free open access medical education in emergency medicine) and the #ChileEM initiative that brings together the programs of the Universidad San Sebastián / MUE, Universidad Católica de Chile and Universidad de Chile, intend to hold joint clinical meetings between the leading training programs, regularly and open to all the health team working in the field of urgency. The specialists already trained are grouped in the Chilean Society of Emergency Medicine (SOCHIMU).
Canada
The two routes to emergency medicine certification can be summarized as follows:
A five-year residency leads to the designation of FRCP(EM) through the Royal College of Physicians and Surgeons of Canada (Emergency Medicine Board Certification – emergency medicine consultant).
A one-year emergency medicine enhanced skills program following a two-year family medicine residency leading to the designation of CCFP(EM) through the College of Family Physicians of Canada (Advanced Competency Certification). The CFPC also allows those are having worked a minimum of four years at a minimum of 400 hours per year in emergency medicine to challenge the examination of special competence in emergency medicine and thus become specialized.
CCFP(EM) emergency physicians outnumber FRCP(EM) physicians by a ratio of about 3 to 1, and they tend to work primarily as clinicians with a minor focus on academic activities such as teaching and research. FRCP(EM) Emergency Medicine Board specialists tend to congregate in academic centres and have more academically oriented careers, which emphasize administration, research, critical care, disaster medicine, and teaching. They also tend to sub-specialize in toxicology, critical care, pediatric emergency medicine, and sports medicine. Furthermore, the FRCP(EM) residency length allows more time for formal training in these areas.
Physician assistants are currently practising in the field of emergency medicine in Canada.
China
The current post-graduate emergency medicine training process is highly complex in China. The first EM post-graduate training took place in 1984 at the Peking Union Medical College Hospital. Because specialty certification in EM has not been established, formal training is not required to practice emergency medicine in China.
About a decade ago, emergency medicine residency training was centralized at the municipal levels, following the Ministry of Public Health guidelines. Residency programs in all hospitals are called residency training bases, which have to be approved by local health governments. These bases are hospital-based, but the residents are selected and managed by the municipal associations of medical education. These associations are also the authoritative body of setting up their residents' training curriculum. All medical school graduates who want to practice medicine have to undergo five years of residency training at designated training bases, the first three years of general rotation followed by two more years of specialty-centred training.
Germany
In Germany, emergency medicine is not handled as a specialization (Facharztrichtung), but any licensed physician can acquire an additional qualification in emergency medicine through an 80-hour course monitored by the respective "Ärztekammer" (medical association, responsible for licensing of physicians). Service as an emergency physician in an ambulance service is part of the specialization training of anaesthesiology. Emergency physicians usually work on a volunteering basis and are often anesthesiologists but maybe specialists of any kind. Especially there is a specialization training in pediatric intensive care.
India
India is an example of how family medicine can be a foundation for emergency medicine training. Many private hospitals and institutes have been providing emergency medicine training for doctors, nurses and paramedics since 1994, with certification programs varying from six months to three years. However, emergency medicine was only recognized as a separate specialty by the Medical Council of India in July 2009.
Malaysia
There are three universities (Universiti Sains Malaysia, Universiti Kebangsaan Malaysia, and Universiti Malaya) that offer master's degrees in emergency medicine – postgraduate training programs of four years in duration with clinical rotations, examinations and a dissertation. The first cohort of locally trained emergency physicians graduated in 2002.
Saudi Arabia
In Saudi Arabia, the Certification of Emergency Medicine takes the four-year Saudi Board of Emergency Medicine (SBEM), which the Saudi Council accredits for Health Specialties (SCFHS). It requires passing the two-part exam: first and final part (written and oral) to obtain the SBEM certificate, equivalent to a doctorate.
Switzerland
Emergency medicine is still not recognised as a fully-fledged specialty in a country that has only recently grasped the importance of having an organised acute medical specialty (during the COVID-19 outbreak). Many attempts to organize the specialty have resulted in a subspecialists training pathway, but to this day, internal medicine, anesthesiology and surgery are still vocally opposed to an emergency medicine specialist title.
United States
Most programs are three years in duration, but some programs are four years long. There are several combined residencies offered with other programs, including family medicine, internal medicine and paediatrics. The US is well known for its excellence in emergency medicine residency programs, leading to some controversy about specialty certification.
There are three ways to become board-certified in emergency medicine:
The American Board of Emergency Medicine (ABEM) is for those with either Doctor of Medicine (MD) or Doctor of Osteopathic Medicine (DO) degrees. The ABEM is under the authority of the American Board of Medical Specialties.
The American Osteopathic Board of Emergency Medicine (AOBEM) certifies only emergency physicians with a DO degree. It is under the authority of the American Osteopathic Association Bureau of Osteopathic Specialists.
The Board of Certification in Emergency Medicine (BCEM) grants board certification in emergency medicine to physicians who have not completed an emergency medicine residency but have completed a residency in other fields (internists, family practitioners, paediatricians, general surgeons, and anesthesiologists). The BCEM is under the authority of the American Board of Physician Specialties.
Several ABMS fellowships are available for emergency medicine graduates, including pre-hospital medicine (emergency medical services), international medicine, advanced resuscitation, hospice and palliative care, research, undersea and hyperbaric medicine, sports medicine, pain medicine, ultrasound, pediatric emergency medicine, disaster medicine, wilderness medicine, toxicology, and critical care medicine.
In recent years, workforce data has led to a recognition of the need for additional training for primary care physicians who provide emergency care. It has led to several supplemental training programs in first-hour emergency care and a few fellowships for family physicians in emergency medicine., and few fellowships for family physicians in emergency medicine.
Funding for training
"In 2010, there were 157 allopathic and 37 osteopathic emergency medicine residency programs, which collectively accept about 2,000 new residents each year. Studies have shown that attending emergency physician supervision of residents correlates to higher quality and more cost-effective practice, primarily when an emergency medicine residency exists." Medical education is primarily funded through the Medicare program; payments are given to hospitals for each resident. "Fifty-five per cent of ED payments come from Medicare, fifteen per cent from Medicaid, five per cent from private payment and twenty-five per cent from commercially insured patients." However, choices of physician specialties are not mandated by any agency or program, so even though emergency departments see many Medicare/Medicaid patients and thus receive much funding for training from these programs, there is still concern over a shortage of specialty-trained emergency medicine providers.
United Kingdom
In the United Kingdom, the Royal College of Emergency Medicine has a role in setting professional standards and assessing trainees. Emergency medical trainees enter specialty training after five or six years of Medical school followed by two years of foundation training. Specialty training takes six years to complete, and success in the assessments and a set of five examinations results in the award of Fellowship of the Royal College of Emergency Medicine (FRCEM).
Historically, emergency specialists were drawn from anaesthesia, medicine, and surgery. Many established EM consultants were surgically trained; some hold the fellowship of Royal College of Surgeons of Edinburgh in accident and emergency – FRCSEd (A&E). trainees in emergency medicine may dual accredit in intensive care medicine or seek sub-specialisation in paediatric emergency medicine.
Turkey
Emergency medicine residencies last four years in Turkey. These physicians have a two-year obligatory service in Turkey to be qualified to have their diploma. After this period, EM specialists can choose to work in private or governmental emergency departments.
Pakistan
The College of Physicians and Surgeons Pakistan accredited the training in emergency medicine in 2010. Emergency medicine training in Pakistan lasts for five years. The initial two years involve trainees being sent to three major areas: medicine and allied, surgery, and allied and critical care. It is divided into six months each, and the rest six months out of the first two years are spent in the emergency department. In the last three years, trainee residents spend most of their time in the emergency room as senior residents. Certificate courses include ACLS, PALS, ATLS, and research and dissertations are required to complete the training successfully. At the end of five years, candidates become eligible to sit for the FCPS part II exam. After fulfilling the requirement, they become fellows of the College of Physicians and Surgeons Pakistan in emergency medicine ().
Institutions providing this training include Shifa International Hospitals Islamabad, Aga Khan University Hospital Karachi, POF Hospital Wah, Lady Reading Hospital Peshawar, Indus Hospital Karachi and Jinnah Post Graduate Medical Center Karachi, and Mayo Hospital, Lahore.
Iran
The first residency program in Iran started in 2002 at Iran University of Medical Sciences, and there are now three-year standard residency programs running in Tehran, Tabriz, Mashhad, Isfahan, and some other universities. All these programs work under the supervision of the emergency medicine specialty board committee. There are now more than 200 (and increasing) board-certified Emergency Physicians in Iran.
Ethical and medicolegal issues
Ethical and medico-legal issues are embedded within the nature of emergency medicine. Issues surrounding competence, end of life care, and the right to refuse care are encountered daily within the emergency department. Of growing significance are the ethical issues and legal obligations that surround the Mental Health Act, as increasing numbers of suicide attempts and self-harm are seen in the emergency department. The Wooltorton case of 2007, in which a patient arrived at the emergency department post overdose with a note specifying her request for no interventions, highlights the dichotomy that often exists between a physician's ethical obligation to "do no harm" and the legality of a patient's right to refuse.
| Biology and health sciences | Fields of medicine | Health |
52975 | https://en.wikipedia.org/wiki/Emergency%20medical%20services | Emergency medical services | Emergency medical services (EMS), also known as ambulance services, pre-hospital care or paramedic services, are emergency services that provide urgent pre-hospital treatment and stabilisation for serious illness and injuries and transport to definitive care. They may also be known as a first aid squad, FAST squad, emergency squad, ambulance squad, ambulance corps, life squad or by other initialisms such as EMAS or EMARS.
In most places, EMS can be summoned by members of the public (as well as medical facilities, other emergency services, businesses and authorities) via an emergency telephone number (such as 911 in the United States) which puts them in contact with a dispatching centre, which will then dispatch suitable resources for the call. Ambulances are the primary vehicles for delivering EMS, though squad cars, motorcycles, aircraft, boats, fire apparatus, and others may be used. EMS agencies may also operate a non-emergency patient transport service, and some have rescue squads to provide technical rescue or search and rescue services.
When EMS is dispatched, they will initiate medical care upon arrival on scene. If it is deemed necessary or a patient requests transport, the unit is then tasked with transferring the patient to the next point of care, typically an emergency department of a hospital. Historically, ambulances only transported patients to care, and this remains the case in parts of the developing world. The term "emergency medical service" was popularised when these services began to emphasise emergency treatment at the scene. In some countries, a substantial portion of EMS calls do not result in a patient being taken to hospital.
Training and qualification levels for members and employees of emergency medical services vary widely throughout the world. In some systems, members may be present who are qualified only to drive ambulances, with no medical training. In contrast, most systems have personnel who retain at least basic first aid certifications, such as basic life support (BLS). In English-speaking countries, they are known as emergency medical technicians (EMTs) and paramedics, with the latter having additional training such as advanced life support (ALS) skills. Physicians and nurses may also provide pre-hospital care to varying degrees in certain countries, a model which is popular in Europe.
History
Precursors
Emergency care in the field has been rendered in different forms since the beginning of recorded history. The New Testament contains the parable of the Good Samaritan, in which a man who has been beaten is cared for by a passing Samaritan. Luke 10:34 (NIV) – "He went to him and bandaged his wounds, pouring on oil and wine. Then he put the man on his own donkey, took him to an inn and took care of him." During the Middle Ages, the Knights Hospitaller were known for rendering assistance to wounded soldiers in the battlefield.
The first use of the ambulance as a specialized vehicle, in battle came about with the ambulances volantes designed by Dominique Jean Larrey (1766–1842), Napoleon Bonaparte's chief surgeon. Larrey was present at the battle of Spires, between the French and Prussians, and was distressed by the fact that wounded soldiers were not picked up by the numerous ambulances (which Napoleon required to be stationed two and half miles back from the scene of battle) until after hostilities had ceased, and set about developing a new ambulance system. Having decided against using the Norman system of horse litters, he settled on two- or four-wheeled horse-drawn wagons, which were used to transport fallen soldiers from the (active) battlefield after they had received early treatment in the field. Larrey's projects for 'flying ambulances' were first approved by the Committee of Public Safety in 1794. Larrey subsequently entered Napoleon's service during the Italian campaigns in 1796, where his ambulances were used for the first time at Udine, Padua and Milan, and he adapted his ambulances to the conditions, even developing a litter which could be carried by a camel for a campaign in Egypt.
Early civilian ambulances
A major advance was made (which in future years would come to shape policy on hospitals and ambulances) with the introduction of a transport carriage for cholera patients in London during 1832. The statement on the carriage, as printed in The Times, said "The curative process commences the instant the patient is put in to the carriage; time is saved which can be given to the care of the patient; the patient may be driven to the hospital so speedily that the hospitals may be less numerous and located at greater distances from each other". This tenet of ambulances providing instant care, allowing hospitals to be spaced further apart, displays itself in modern emergency medical planning.
The first known hospital-based ambulance service operated out of Commercial Hospital, Cincinnati, Ohio (now the Cincinnati General) by 1865. This was soon followed by other services, notably the New York service provided out of Bellevue Hospital which started in 1869 with ambulances carrying medical equipment, such as splints, a stomach pump, morphine, and brandy, reflecting contemporary medicine.
Another early ambulance service was founded by Jaromir V. Mundy, Count J. N. Wilczek, and Eduard Lamezan-Salins in Vienna after the disastrous fire at the Vienna Ringtheater in 1881. Named the "Vienna Voluntary Rescue Society," it served as a model for similar societies worldwide.
In June 1887 the St John Ambulance Brigade was established to provide first aid and ambulance services at public events in London. It was modelled on a military-style command and discipline structure.
Motorization
Also in the late 19th century, the automobile was being developed, and in addition to horse-drawn models, early 20th century ambulances were powered by steam, gasoline, and electricity, reflecting the competing automotive technologies then in existence. However, the first motorized ambulance was brought into service in the last year of the 19th century, with the Michael Reese Hospital, Chicago, taking delivery of the first automobile ambulance, donated by 500 prominent local businessmen, in February 1899. This was followed in 1900 by New York City, who extolled its virtues of greater speed, more safety for the patient, faster stopping and a smoother ride. These first two automobile ambulances were electrically powered with 2 hp motors on the rear axle.
During World War I, further advances were made in providing care before and during transport; traction splints were introduced during the war and were found to have a positive effect on the morbidity and mortality of patients with leg fractures. Two-way radios became available shortly after World War I, enabling for more efficient radio dispatch of ambulances in some areas. Prior to World War II, there were some areas where a modern ambulance carried advanced medical equipment, was staffed by a physician, and was dispatched by radio. In many locations, however, ambulances were hearses, the only available vehicle that could carry a recumbent patient, and were thus frequently run by funeral homes. These vehicles, which could serve either purpose, were known as combination cars.
Prior to World War II, hospitals provided ambulance service in many large cities. With the severe manpower shortages imposed by the war effort, it became difficult for many hospitals to maintain their ambulance operations. City governments in many cases turned ambulance services over to the police or fire department. No laws required minimal training for ambulance personnel and no training programs existed beyond basic first aid. In many fire departments, assignment to ambulance duty became an unofficial form of punishment.
Rise of modern EMS
Advances in the 1960s, especially the development of CPR and defibrillation as the standard form of care for out-of-hospital cardiac arrest, along with new pharmaceuticals, led to changes in the tasks of the ambulances. In Belfast, Northern Ireland the first experimental mobile coronary care ambulance successfully resuscitated patients using these technologies. Freedom House Ambulance Service was the first civilian emergency medical service in the United States to be staffed by paramedics, all of whom were African-American.
One well-known report in the US during that time was Accidental Death and Disability: The Neglected Disease of Modern Society, also known as The White Paper. The report concluded that ambulance services in the US varied widely in quality and were often unregulated and unsatisfactory. These studies placed pressure on governments to improve emergency care in general, including the care provided by ambulance services. The government reports resulted in the creation of standards in ambulance construction concerning the internal height of the patient care area (to allow for an attendant to continue to care for the patient during transport), and the equipment (and thus weight) that an ambulance had to carry, and several other factors.
In 1971 a progress report was published at the annual meeting, by the then president of American Association of Trauma, Sawnie R. Gaston M.D. Dr. Gaston reported the study was a "superb white paper" that "jolted and wakened the entire structure of organized medicine." This report is created as a "prime mover" and made the "single greatest contribution of its kind to the improvement of emergency medical services". Since this time a concerted effort has been undertaken to improve emergency medical care in the pre-hospital setting. Such advancements included Dr. R Adams Cowley creating the country's first statewide EMS program, in Maryland.
The developments were paralleled in other countries. In the United Kingdom, a 1973 law merged the municipal ambulance services into larger agencies and set national standards. In France, the first official SAMU agencies were founded in the 1970s.
Organization
Depending on country, area within country, or clinical need, emergency medical services may be provided by one or more different types of organization. This variation may lead to large differences in levels of care and expected scope of practice. Some countries closely regulate the industry (and may require anyone working on an ambulance to be qualified to a set level), whereas others allow quite wide differences between types of operator.
Municipal "third service" ambulance service
Operating separately from (although alongside) the fire and police services of the area, these ambulances are funded by local, provincial or national governments. In some countries, these only tend to be found in big cities, whereas in countries such as the United Kingdom, almost all emergency ambulances are part of a national health system.
In the United States, ambulance services provided by a local government are often referred to as "third service" EMS (the fire department, police department, and EMS department forming an emergency services trio) by the members of said service, as well as other city officials and residents. The most notable examples of this model in the United States are Pittsburgh Bureau of Emergency Medical Services, Boston EMS, New Orleans Emergency Medical Services, and Cleveland EMS. Government ambulance services also have to take civil service exams just like government fire departments and police. In the United States, certain federal government agencies employ emergency medical technicians at the basic and advanced life support levels, such as the National Park Service and the Federal Bureau of Prisons.
Fire- or police-linked service
In countries such as the United States, Japan, France, South Korea and parts of India, ambulances can be operated by the local fire or police services. Fire-based EMS is the most common model in the United States, where nearly all urban fire departments provide EMS and a majority of emergency transport ambulance services in large cities are part of fire departments. Examples of this model are the New York City Fire Department (FDNY) and the Baltimore City Fire Department.
It is rare for a police department in the United States to provide EMS or ambulance services, although many police officers have basic medical training (such as Nalaxone use and CPR). One notable example is New Orleans Emergency Medical Services, which was formed as a hospital-based service, was operated by the New Orleans Police Department from 1947 to 1985, and is currently operated by the New Orleans Health Department and the New Orleans Office of Homeland Security and Emergency Preparedness, separate from the New Orleans Fire Department.
Charity/not-for-profit ambulance service
Charities or non-profit ambulance departments operate some emergency medical services. They are primarily staffed by volunteers, though many also have paid personnel. These may be linked to a volunteer fire service, and some volunteers may provide both services. Some ambulance charities specialize in providing cover at public gatherings and events (e.g. sporting events), while others provide care to the wider community.
The International Red Cross and Red Crescent Movement is the largest charity in the world that provides emergency medicine. (in some countries, it operates as a private ambulance service). Other organisations include St John Ambulance, the Order of Malta Ambulance Corps and Hatzalah, as well as small local volunteer/paid departments. In the United States, volunteer ambulances are rarer, but can still be seen in both metropolitan and rural areas (e.g. Hatzalah). Charities such as BASICS Scotland, specialise in facilitating training medical professionals to volunteer to assist the statutory ambulance services in the care of patients, through their attendance at those with serious illnesses or injuries.
A few charities provide ambulances for taking patients on trips or vacations away from hospitals, hospices or care homes where they are in long-term care. Examples include the UK's Jumbulance project.
Private/corporate ambulance service
Some ambulances are operated by commercial companies with paid employees, usually on a contract to the local or national government, Hospital Networks, Health Care Facilities and Insurance Companies.
In the U.S., private ambulance companies provide emergency medical services in large cities and rural areas by contracting with local governments. In areas where the local county or city provide their own emergency services, private companies provide discharges and transfers from hospitals and to/from other health related facilities and homes. In most areas private companies are part of the local government emergency disaster plan, and are relied upon for the overall EMS response, treatment and recovery.
In some areas, private companies may provide only the patient transport elements of ambulance care (i.e. non-urgent), but in some places, they are contracted to provide emergency care, or to form a 'second tier' response, where they only respond to emergencies when all of the full-time emergency ambulance crews are busy. This may mean that a government or other service provide the 'emergency' cover, whilst a private firm may be charged with 'minor injuries' such as cuts, bruises or even helping the mobility-impaired if they have for example fallen and simply need help to get up again, but do not need treatment. This system has the benefit of keeping emergency crews available at all times for genuine emergencies. These organisations may also provide services known as 'Stand-by' cover at industrial sites or at special events. In Latin America, private ambulance companies are often the only readily-available EMS service
Combined emergency service
These are full service emergency service agencies, which may be found in places such as airports or large colleges and universities like for example the UCLA EMS.Their key feature is that all personnel are trained not only in ambulance (EMT) care, but as a firefighter and a peace officer (police function). They may be found in smaller towns and cities, where demand or budget is too low to support separate services. This multi-functionality allows to make the most of limited resource or budget, but having a single team respond to any emergency.
Hospital-based service
Hospitals or larger hospital systems may provide their own ambulance service as a service to the community, or where ambulance care is unreliable or chargeable. Many hospital-based EMS departments operate solely with their hospital, though some operate more independently and can transport patients to whichever hospital may be needed or desired.
Internal ambulances
Many large factories and other industrial centers, such as chemical plants, oil refineries, breweries, and distilleries, have emergency medical services provided by employers as a means of protecting their interests and the welfare of their staff. These are often used as first response vehicles in the event of a fire or explosion.
Purpose
Emergency medical services exists to fulfill the basic principles of first aid, which are to Preserve Life, Prevent Further Injury, and Promote Recovery. This common theme in medicine is demonstrated by the "star of life". The Star of Life shown here, where each of the 'arms' to the star represent one of the six points, are used to represent the six stages of high quality pre-hospital care, which are:
Early detection – members of the public, or another agency, find the incident and understand the problem
Early reporting – the first persons on scene make a call to the emergency medical services (911) and provide details to enable a response to be mounted
Early response – the first professional (EMS) rescuers are dispatched and arrive on scene as quickly as possible, enabling care to begin
Good on-scene/field care – the emergency medical service provides appropriate and timely interventions to treat the patient at the scene of the incident without doing further harm.
Care in transit -– the emergency medical service load the patient in to suitable transport and continue to provide appropriate medical care throughout the journey
Transfer to definitive care – the patient is handed over to an appropriate care setting, such as the emergency department at a hospital, in to the care of physicians
Strategies for delivering care
Although a variety of differing philosophical approaches are used in the provision of EMS care around the world, they can generally be placed into one of two categories; one physician-led and the other led by pre-hospital allied health staff such as emergency medical technicians or paramedics. These models are commonly referred to as the Franco-German model and Anglo-American model.
Studies have been inconclusive as to whether one model delivers better results than the other. A 2010 study in the Oman Medical Journal suggested that rapid transport was a better strategy for trauma cases, while stabilization at the scene was a better strategy for cardiac arrests.
Levels of care
Many systems have tiers of response for medical emergencies. For example, a common arrangement in the United States is that fire engines or volunteers are sent to provide a rapid initial response to a medical emergency, while an ambulance is sent to provide advanced treatment and transport the patient. In France, fire service and private company ambulances provide basic care, while hospital-based ambulances with physicians on board provide advanced care. In many countries, an air ambulance provides a higher level of care than a regular ambulance.
Examples of level of care include:
First aid consists of basic skills that are commonly taught to members of the public, such as cardiopulmonary resuscitation, bandaging wounds and saving someone from choking.
Basic Life Support (BLS) is often the lowest level of training that can be held by those who treat patients on an ambulance. Commonly, it includes administering oxygen therapy, some drugs and a few invasive treatments. BLS personnel may either operate a BLS ambulance on their own, or assist a higher qualified crewmate on an ALS ambulance. In English-speaking countries, BLS ambulance crew members are known as emergency medical technicians or emergency care assistants.
Intermediate Life Support (ILS), also known as Limited Advanced Life Support (LALS), is positioned between BLS and ALS but is less common than both. It is commonly a BLS provider with a moderately expanded skill set, but where it is present it usually replaces BLS.
Advanced Life Support (ALS) has a considerably expanded range of skills such as intravenous therapy, cricothyrotomy and interpreting an electrocardiogram. The scope of this higher tier response varies considerably by country. Paramedics commonly provide ALS, but some countries require it to be a higher level of care and instead employ physicians in this role. Additionally Advanced Life Support includes administering therapeutic doses of electrical shock to those who are in cardiac arrest or using drugs to stimulate the heart, Airway therapy, and so on and so forth. Most ambulances are equipped with advanced Life Support equipment and have paramedics on board. While some fire departments have ambulances, first aid and squads utilize ambulances for emergency medical services.
Critical Care Transport (CCT), also known as medical retrieval or rendez vous MICU protocol in some countries (Australia, NZ, Great Britain, and Francophone Canada) refers to the critical care transport of patients between hospitals (as opposed to pre-hospital). Such services are a key element in regionalized systems of hospital care where intensive care services are centralized to a few specialist hospitals. An example of this is the Emergency Medical Retrieval Service in Scotland. This level of care is likely to involve traditional healthcare professionals (in addition to or instead of critical care-trained paramedics), meaning nurses and/or physicians working in the pre-hospital setting and even on ambulances.
Transport-only
The most basic emergency medical services are provided as a transport operation only, simply to take patients from their location to the nearest medical treatment. This was historically the case in all countries. It remains the case in much of the developing world, where operators as diverse as taxi drivers and undertakers may transport people to hospital.
Transport-centered EMS
The Anglo-American model is also known as "load and go" or "scoop and run". In this model, ambulances are staffed by paramedics and/or emergency medical technicians. They have specialized medical training, but not to the same level as a physician. In this model it is rare to find a physician actually working routinely in ambulances, although they may be deployed to major or complex cases. The physicians who work in EMS provide oversight for the work of the ambulance crews. This may include off-line medical control, where they devise protocols or 'standing orders' (procedures for treatment). This may also include on-line medical control, in which the physician is contacted via radio or phone to provide advice and authorization for various medical interventions or for a patient's desire to refuse care.
In some cases, such as in the UK, South Africa and Australia, a paramedic may be an autonomous medical professional, and does not require the permission of a physician to administer interventions or medications from an agreed list, and can perform roles such as suturing or prescribing medication to the patient. Recently "Telemedicine" has been making an appearance in ambulances. Similar to online medical control, this practice allows paramedics to remotely transmit data such as vital signs and 12 and 15 lead ECGs to the hospital from the field. This allows the emergency department to prepare to treat patients prior to their arrival. This is allowing lower level providers (Such as EMT-B) in the United States to utilize these advanced technologies and have the doctor interpret them, thus bringing rapid identification of rhythms to areas where paramedics are stretched thin. While most insurance companies only reimburse EMS providers for transporting patients to 911 receiving facilities (e.g. Emergency Departments),the Center to Medicare and Medicaid Services is in the process of evaluating a payment model to enable reimbursement for patients evaluated and treated on-scene.
Major trauma
The essential decision in prehospital care is whether the patient should be immediately taken to the hospital, or advanced care resources are taken to the patient where they lie. The "scoop and run" approach is exemplified by the MEDEVAC aeromedical evacuation helicopter, whereas the "stay and play" is exemplified by the French and Belgian SMUR emergency mobile resuscitation unit or the German "Notarzt"-System (preclinical emergency physician).
The strategy developed for prehospital trauma care in North America is based on the Golden Hour theory, i.e., that a trauma victim's best chance for survival is in an operating room, with the goal of having the patient in surgery within an hour of the traumatic event. This appears to be true in cases of internal bleeding, especially penetrating trauma such as gunshot or stab wounds. Thus, minimal time is spent providing prehospital care (spine immobilization; "ABCs", i.e. ensure airway, breathing and circulation; external bleeding control; endotracheal intubation) and the victim is transported as fast as possible to a trauma centre.
The aim in "Scoop and Run" treatment is generally to transport the patient within ten minutes of arrival, hence the birth of the phrase, "the platinum ten minutes" (in addition to the "golden hour"), now commonly used in EMT training programs. The "Scoop and Run" is a method developed to deal with trauma, rather than strictly medical situations (e.g. cardiac or respiratory emergencies), however, this may be changing. Increasingly, research into the management of S-T segment elevation myocardial infarctions (STEMI) occurring outside of the hospital, or even inside community hospitals without their own PCI labs, suggests that time to treatment is a clinically significant factor in heart attacks, and that trauma patients may not be the only patients for whom 'load and go' is clinically appropriate. In such conditions, the gold standard is the door to balloon time. The longer the time interval, the greater the damage to the myocardium, and the poorer the long-term prognosis for the patient. Current research in Canada has suggested that door to balloon times are significantly lower when appropriate patients are identified by paramedics in the field, instead of the emergency room, and then transported directly to a waiting PCI lab. The STEMI program has reduced STEMI deaths in the Ottawa region by 50 per cent. In a related program in Toronto, EMS has begun to use a procedure of 'rescuing' STEMI patients from the Emergency Rooms of hospitals without PCI labs, and transporting them, on an emergency basis, to waiting PCI labs in other hospitals.
Physician-led EMS
Physician-led EMS is also known as the Franco-German model, "stay and play", "stay and stabilize" or "delay and treat". In a physician-led system, doctors respond directly to all major emergencies requiring more than simple first aid. The physicians will attempt to treat casualties at the scene and will only transport them to hospital if it is deemed necessary. If patients are transported to hospital, they are more likely to go straight to a ward rather than to an emergency department. Countries that use this model include Austria, France, Belgium, Luxembourg, Italy, Spain, Brazil and Chile.
In some cases in this model, such as France, there is no direct equivalent to a paramedic. Physicians and (in some cases) nurses provide all medical interventions for the patient. Other ambulance personnel are not non-medically trained and only provide driving and heavy lifting. In other applications of this model, as in Germany, a paramedic equivalent does exist, but is an assistant to the physician with a restricted scope of practice. They are only permitted to perform Advanced Life Support (ALS) procedures if authorized by the physician, or in cases of immediate life-threatening conditions. Ambulances in this model tend to be better equipped with more advanced medical devices, in essence, bringing the emergency department to the patient. High-speed transport to hospitals is considered, in most cases, to be unnecessarily unsafe, and the preference is to remain and provide definitive care to the patient until they are medically stable, and then accomplish transport. In this model, the physician and nurse may actually staff an ambulance along with a driver, or may staff a rapid response vehicle instead of an ambulance, providing medical support to multiple ambulances.
Personnel
Ambulance personnel are generally professionals and in some countries their use is controlled through training and registration. While these job titles are protected by legislation in some countries, this protection is by no means universal, and anyone might, for example, call themselves an 'EMT' or a 'paramedic', regardless of their training, or the lack of it. In some jurisdictions, both technicians and paramedics may be further defined by the environment in which they operate, including such designations as 'Wilderness', 'Tactical', and so on.
A unique aspect of EMS is that there are two hierarchies of authority, as the chain of command is separate to medical authority.
Basic life support (BLS)
Emergency medical dispatcher
An emergency medical dispatcher is also called an EMD. An increasingly common addition to the EMS system is the use of highly trained dispatch personnel who can provide "pre-arrival" instructions to callers reporting medical emergencies. They use carefully structured questioning techniques and provide scripted instructions to allow callers or bystanders to begin definitive care for such critical problems as airway obstructions, bleeding, childbirth, and cardiac arrest. Even with a fast response time by a first responder measured in minutes, some medical emergencies evolve in seconds. Such a system provides, in essence, a "zero response time," and can have an enormous impact on positive patient outcomes.
First responder
Certified first responders may be sent to provide first aid, sometimes to an advanced level. Their duties include the provision of immediate life-saving care in the event of a medical emergency; commonly advanced first aid, oxygen administration, cardio-pulmonary resuscitation (CPR), and automated external defibrillator (AED) usage. The first responder training is considered a bare minimum for emergency service workers who may be sent out in response to an emergency call. First responders are commonly dispatched by the ambulance service to arrive quickly and stabilize the patient before the ambulance arrives, and to then assist the ambulance crew.
Some EMS agencies have set up volunteer schemes, who can be dispatched to a medical emergency before the ambulance arrives. Examples of this include Community First Responder schemes run by ambulance services the UK and similar volunteer schemes operated by the fire services in France. In some countries such as the US, there may be autonomous groups of volunteer responders such as rescue squads. Police officers and firefighters who are on duty for another emergency service may also be deployed in this role, though some firefighters are trained to a more advanced medical level.
Besides first responders who are deployed to an emergency, there are others who may be stationed at public events. The International Red Cross and Red Crescent Movement and St John Ambulance both provide first aiders in these roles.
Ambulance driver
Some agencies separate the 'driver' and 'attendant' functions, employing ambulance driving staff with no medical qualification (or just a first aid and CPR certificates), whose job is to drive ambulances. While this approach persists in some countries, such as India, it is generally becoming increasingly rare. Ambulance drivers may be trained in radio communications, ambulance operations and emergency response driving skills.
Non-emergency attendant
Many countries employ ambulance staff who only carry out non-emergency patient transport duties (which can include stretcher or wheelchair cases). Dependent on the provider (and resources available), they may be trained in first aid or extended skills such as use of an AED, oxygen therapy, pain relief and other live-saving or palliative skills. In some services, they may also provide emergency cover when other units are not available, or when accompanied by a fully qualified technician or paramedic. The role is known as an Ambulance Care Assistant in the United Kingdom.
Emergency care assistant
Emergency care assistants are of a frontline under both emergency and non-emergency conditions to incidents. Their role is to assist the clinician that they are working with, either a Technician or Paramedic, in their duties, whether that be drawing up drugs, setting up fluids (but not attaching), doing basic observations or performing 12 lead ECG assessments.
Emergency medical technician
Emergency medical technicians are usually able to perform a wide range of emergency care skills, such as automated defibrillation, care of spinal injuries and oxygen therapy. In few jurisdictions, some EMTs are able to perform duties as IV and IO cannulation, administration of a limited number of drugs (including but not limited to Epinephrine, Narcan, Oxygen, Aspirin, Nitroglycerin – dependent on country, state, and medical direction), more advanced airway procedures, CPAP, and limited cardiac monitoring. Most advanced procedures and skills are not within the national scope of practice for an EMT. As such most states require additional training and certifications to perform above the national curriculum standards. In the United States, an EMT certification requires intense courses and training in field skills. A certification expires after two years and holds a requirement of taking 48 CEUs (continuing education credits). 24 of these credits must be in refresher courses while the other 24 can be taken in a variety ways such as emergency driving training, pediatric, geriatric, or bariatric care, specific traumas, etc. Is usually made up of 3 levels in the US. EMT-B, EMT-I (EMT-A in some states) and EMT-Paramedic. The National Registry of EMT New Educational Standards for EMS renamed the provider levels as follows: Emergency Medical Responder (EMR), Emergency Medical Technician (EMT-B), Advanced EMT (AEMT), and Paramedic (EMT-P).
Advanced life support (ALS)
Paramedic
A paramedic has a high level of pre-hospital medical training and usually involves key skills not performed by technicians, often including cannulation (and with it the ability to use a range of drugs to relieve pain, correct cardiac problems, and perform endotracheal intubation), cardiac monitoring, 12-lead ECG interpretation, ultrasound, intubation, pericardiocentesis, cardioversion, thoracostomy, and other skills such as performing a surgical cricothyrotomy. The most important function of the paramedic is to identify and treat any life-threatening conditions and then to assess the patient carefully for other complaints or findings that may require emergency treatment. In many countries, this is a protected title, and use of it without the relevant qualification may result in criminal prosecution. In the United States, paramedics represent the highest licensure level of prehospital emergency care. In addition, several certifications exist for Paramedics such as Wilderness ALS Care, Flight Paramedic Certification (FP-C), and Critical Care Emergency Medical Transport Program certification.
Critical care paramedic
A critical care paramedic, also known as an advanced practice paramedic or specialist paramedic, is a paramedic with additional training to deal with critically ill patients. Critical care paramedics often work on air ambulances, which are more likely to be dispatched to emergencies requiring advanced care skills. They may also work on land ambulances. The training, permitted skills, and certification requirements vary from one jurisdiction to the next. It also varies to whether they are trained externally by a university or professional body or 'in house' by their EMS agency.
These providers have a vast array of and medications to handle complex medical and trauma patients. Examples of medication are dobutamine, dopamine, Propofol, blood and blood products to name just a few. Some examples of skills include, but not limited to, life support systems normally restricted to the ICU or critical care hospital setting such as mechanical ventilators, Intra-aortic balloon pump (IABP) and external pacemaker monitoring. Depending on the service medical direction, these providers are trained on placement and use of UVCs (Umbilical Venous Catheter), UACs (Umbilical Arterial Catheter), surgical airways, central lines, arterial lines and chest tubes.
Emergency care practitioner
In the United Kingdom and South Africa, some serving paramedics receive additional university education to become practitioners in their own right, which gives them absolute responsibility for their clinical judgement, including the ability to autonomously prescribe medications. An emergency care practitioner or paramedic practitioner is a position that is designed to bridge the link between ambulance care and the care of a general practitioner. ECPs are university graduates in Emergency Medical Care or qualified paramedics who have undergone further training, and are authorized to perform specialized techniques. Additionally some may prescribe medicines for longer-term care, such as antibiotics and in the United Kingdom they can prescribe a broad range of medicines. With respect to a Primary Health Care setting, they are also educated in a range of Diagnostic techniques.
Traditional healthcare professionals in EMS
Registered nurses
The use of registered nurses (RNs) in the pre-hospital setting is common in many countries in absence of paramedics. In some regions of the world nurses are the primary healthcare worker that provides emergency medical services. In European countries such as France or Italy, also use nurses as a means of providing ALS services. These nurses may work under the direct supervision of a physician, or, in rarer cases, independently. In some places in Europe, notably Norway, paramedics do exist, but the role of the 'ambulance nurse' continues to be developed, as it is felt that nurses may bring unique skills to some situations encountered by ambulance crews.
In North America, and to a lesser extent elsewhere in the English-speaking world, some jurisdictions use specially trained nurses for medical transport work. These are mostly air-medical personnel or critical care transport providers, often working in conjunction with a technician, paramedic or physician on emergency interfacility transports. In the United States, the most common uses of ambulance-based registered nurses is in the Critical Care/Mobile Intensive Care transport, and in Aeromedical EMS. Such nurses are normally required by their employers (in the US) to seek additional certifications beyond the primary nursing licensure. Four individual states have an Intensive Care or Prehospital Nurse licensure. Many states allow registered nurses to also become registered paramedics according to their role in the emergency medical services team. In Estonia 60% of ambulance teams are led by nurse. Ambulance nurses can do almost all emergency procedures and administer medicines pre-hospital such as physicians in Estonia. In the Netherlands, all ambulances are staffed by a registered nurse with additional training in emergency nursing, anesthesia or critical care, and a driver-EMT. In Sweden, since 2005, all emergency ambulances should be staffed by at least one registered nurse since only nurses are allowed to administer drugs. And all Advanced Life Support Ambulances are staffed at least by a registered nurse in Spain. In France, since 1986, fire department-based rescue ambulances have had the option of providing resuscitation service (reanimation) using specially trained nurses, operating on protocols, while SAMU-SMUR units are staffed by physicians and nurses
Physician
In countries with a physician-led EMS model, such as France, Italy, the German-speaking countries (Germany, Switzerland, Austria), and Spain, physicians respond to all cases that require more than basic first aid. In some versions of this model (such as France, Italy, and Spain), there is no direct equivalent to a paramedic, as ALS is performed by physicians. In the German-speaking countries, paramedics are assistants to ambulance physicians (called Notarzt). In these countries, if a physician is present, paramedics require permission from the physician to administer treatments such as defibrillation and drugs. If there is no physician on scene and a life-threatening condition is present, they may administer treatments that follow the physician's instructions.
In countries where EMS is led by paramedics, the ambulance service may still employ physicians. They may serve on specialist response vehicles, such as the air ambulances in the UK. They may also provide advice and devise protocols for treatment, with a medical director acting as the most senior medical adviser to the ambulance service. In the United States, EMS became an officially recognized subspecialty by the American Board of Emergency Medicine in 2010, and the first examinations were held in 2013. Many states now recommend EMS board certification for newly hired medical directors of EMS agencies.
Specialist EMS
Air ambulances
Air ambulances (also known as medevac) often complement a land ambulance service. In some remote areas, they may even form the primary ambulance service. Like many innovations in EMS, medical aircraft were first used in the military. One of the first recorded aircraft rescues of a casualty was in 1917 in Turkey, when a soldier in the Camel Corps who had been shot in the ankle was flown to hospital in a de Havilland DH9. In 1928, the first civilian air medical service was founded in Australia to provide healthcare to people living in remote parts of the Outback. This service became the Royal Flying Doctor Service. The use of helicopters was pioneered in the Korean War, when time to reach a medical facility was reduced from 8 hours to 3 hours in World War II, and again to 2 hours by the Vietnam War.
Aircraft can travel faster and operate in a wider coverage area than a land ambulance. They have a particular advantage for major trauma injuries, especially when they occur in rural or isolated areas. The well-established theory of the golden hour suggests that major trauma patients should be transported as quickly as possible to a specialist trauma center. Therefore, medical first responders in a helicopter can provide both a higher level of care at the scene, faster transport to a specialist hospital and critical care during the journey. A disadvantage is that it can be dangerous and potentially not possible for them to fly at night or in bad weather.
Tactical (hazardous area)
Some EMS agencies have set up specialist teams to help those injured in a major incident or a dangerous situation. These include tactical police operations, active shooters, bombings, hazmat situations, building collapses, fires and natural disasters. In the US, these are often known as Tactical EMS teams and are often deployed alongside police SWAT teams. The equivalent in UK ambulance services is a Hazardous Area Response Team (HART).
Wilderness
Wilderness EMS-like systems (WEMS) have been developed to provide medical responses in remote areas, which may have significantly different needs to an urban area. Examples include the National Ski Patrol or the regional-responding Appalachian Search and Rescue Conference (USA based). Like traditional EMS providers, all wilderness emergency medical (WEM) providers must still operate under on-line or off-line medical oversight. To assist physicians in the skills necessary to provide this oversight, the Wilderness Medical Society and the National Association of EMS Physicians jointly supported the development in 2011 of a unique "Wilderness EMS Medical Director" certification course, which was cited by the Journal of EMS as one of the Top 10 EMS Innovations of 2011. Skills taught in WEMT courses exceeding the EMT-Basic scope of practice include catheterization, antibiotic administration, use of intermediate Blind Insertion Airway Devices (i.e. King Laryngeal Tube), Nasogastric Intubation, and simple suturing; however, the scope of practice for the WEMT still falls under BLS level care. A multitude of organizations provide WEM training, including private schools, non-profit organizations such as the Appalachian Center for Wilderness Medicine and the Wilderness EMS Institute, military branches, community colleges and universities, EMS-college-hospital collaborations, and others.
Occupational Health Hazards in EMS
Hazards
The rate of occupational injuries for EMS workers is greater than the rate for the general population. Occupational hazards for health professionals are well studied and generally apply to EMS. Occupational health hazards in emergency medical services include lifting injuries, violent patients, transportation incidents and harmful exposures (exposure to loud noises, chemicals and infectious diseases).
The National Institute of Occupational Safety and Health (NIOSH) established that the most common injury/illness of EMS clinicians are sprains/strains (41%) and exposures (20%). The occupational hazard with the greatest rate of lethality in Emergency Medical Services is ground and air ambulance crashes. Crashes are uncommon and Injuries from crashes make up less than 8% of injuries to EMS clinicians.
Hazard Mitigation
In the US, OSHA and CDC associated NIOSH have published guidelines to protect EMS workers from the occupational hazards presented by the job's requirements. These guidelines include sanitization practices, PPE requirements, and fatigue protocols. Fundamental facets of safety, such as providing workers with clear roles and clear communication of expectations to discourage risk-taking and encourage safe practices are essential to minimize occupational hazards.
Overexertion injuries can be avoided with core training, increasing flexibility, and improving muscular endurance. Safe lifting form is encouraged. Exposure to chemical, biological, sensory and physical hazards can all be mitigated with an increased use of PPE. Psychosocial hazards such as job stress and exposure to violence or trauma can be managed with peer support resources for emergency responders who are suffering mental health effects. Most (80%) of states have established independent agency critical incident stress debriefing/management (CISD/M) programs that come into action whenever an incident is deemed particularly traumatic. Transportation related injury and fatality can be minimized with better driver training programs and protocols.
Organization in different countries
Emergency medical services in Austria
Emergency medical services in Australia
Emergency medical services in Belgium
Emergency medical services in Brazil
Emergency medical services in Canada
Emergency medical services in Finland
Emergency medical services in France
Emergency medical services in Germany
Emergency medical services in Hong Kong
Emergency medical services in Iceland
Emergency medical services in Iran
Emergency medical services in Ireland
Emergency medical services in Israel
Emergency medical services in Italy
Emergency medical services in the Netherlands
Emergency medical services in New Zealand
Emergency medical services in Norway
Emergency medical services in Pakistan
Emergency medical services in Poland
Emergency medical services in Portugal
Emergency medical services in Romania
Emergency medical services in Russia
Emergency medical services in Spain
Emergency medical services in South Africa
Emergency medical services in Sri Lanka
Emergency medical services in Ukraine
Emergency medical services in the United Kingdom
Emergency medical services in the United States
| Biology and health sciences | General concepts | null |
52999 | https://en.wikipedia.org/wiki/Phobos%20%28moon%29 | Phobos (moon) | Phobos (; systematic designation: ) is the innermost and larger of the two natural satellites of Mars, the other being Deimos. The two moons were discovered in 1877 by American astronomer Asaph Hall. Phobos is named after the Greek god of fear and panic, who is the son of Ares (Mars) and twin brother of Deimos.
Phobos is a small, irregularly shaped object with a mean radius of . It orbits from the Martian surface, closer to its primary body than any other known natural satellite to a planet. It orbits Mars much faster than Mars rotates and completes an orbit in just 7 hours and 39 minutes. As a result, from the surface of Mars it appears to rise in the west, move across the sky in 4 hours and 15 minutes or less, and set in the east, twice each Martian day. Phobos is one of the least reflective bodies in the Solar System, with an albedo of 0.071. Surface temperatures range from about on the sunlit side to on the shadowed side. The notable surface feature is the large impact crater, Stickney, which takes up a substantial proportion of the moon's surface. The surface is also marked by many grooves, and there are numerous theories as to how these grooves were formed.
Images and models indicate that Phobos may be a rubble pile held together by a thin crust that is being torn apart by tidal interactions. Phobos gets closer to Mars by about per year.
Discovery
Phobos was discovered by the American astronomer Asaph Hall on 18 August 1877 at the United States Naval Observatory in Washington, D.C., at about 09:14 Greenwich Mean Time. (Contemporary sources, using the pre-1925 astronomical convention that began the day at noon, give the time of discovery as 17 August at 16:06 Washington mean time, meaning 18 August 04:06 in the modern convention.) Hall had discovered Deimos, Mars' other moon, a few days earlier.
The discoveries were made using the world's largest refracting telescope, the 26-inch "Great Equatorial".
The names, originally spelled Phobus and Deimus respectively, were suggested by the British academic Henry Madan, a science master at Eton College, who based them on Greek mythology, in which Phobos is a companion to the god, Ares.
Physical characteristics
Phobos has dimensions of , and retains too little mass to be rounded under its own gravity. Phobos does not have an atmosphere due to its low mass and low gravity. It is one of the least reflective bodies in the Solar System, with an albedo of about 0.071. Infrared spectra show that it has carbon-rich material found in carbonaceous chondrites, and its composition shows similarities to that of Mars' surface. Phobos' density is too low to be solid rock, and it is known to have significant porosity. These results led to the suggestion that Phobos might contain a substantial reservoir of ice. Spectral observations indicate that the surface regolith layer lacks hydration, but ice below the regolith is not ruled out. Surface temperatures range from about on the sunlit side to on the shadowed side.
Unlike Deimos, Phobos is heavily cratered, with one of the craters near the equator having a central peak despite the moon's small size. The most prominent of these is the crater Stickney, a large impact crater some in diameter, which takes up a substantial proportion of the moon's surface area. As with Mimas' crater Herschel, the impact that created Stickney must have nearly shattered Phobos.
Many grooves and streaks also cover the oddly shaped surface. The grooves are typically less than deep, wide, and up to in length, and were originally assumed to have been the result of the same impact that created Stickney. Analysis of results from the Mars Express spacecraft, however, revealed that the grooves are not radial to Stickney, but are centered on the leading apex of Phobos in its orbit (which is not far from Stickney). Researchers suspect that they have been excavated by material ejected into space by impacts on the surface of Mars. The grooves thus formed as crater chains, and all of them fade away as the trailing apex of Phobos is approached. They have been grouped into 12 or more families of varying age, presumably representing at least 12 Martian impact events. However, in November 2018, following further computational probability analysis, astronomers concluded that the many grooves on Phobos were caused by boulders, ejected from the asteroid impact that created Stickney crater. These boulders rolled in a predictable pattern on the surface of the moon.
Faint dust rings produced by Phobos and Deimos have long been predicted but attempts to observe these rings have, to date, failed. Recent images from Mars Global Surveyor indicate that Phobos is covered with a layer of fine-grained regolith at least 100 meters thick; it is hypothesized to have been created by impacts from other bodies, but it is not known how the material stuck to an object with almost no gravity.
The unique Kaidun meteorite that fell on a Soviet military base in Yemen in 1980 has been hypothesized to be a piece of Phobos, but this couldn't be verified because little is known about the exact composition of Phobos.
Shklovsky's "Hollow Phobos" hypothesis
In the late 1950s and 1960s, the unusual orbital characteristics of Phobos led to speculations that it might be hollow. Around 1958, Russian astrophysicist Iosif Samuilovich Shklovsky, studying the secular acceleration of Phobos' orbital motion, suggested a "thin sheet metal" structure for Phobos, a suggestion which led to speculations that Phobos was of artificial origin. Shklovsky based his analysis on estimates of the upper Martian atmosphere's density, and deduced that for the weak braking effect to be able to account for the secular acceleration, Phobos had to be very light—one calculation yielded a hollow iron sphere across but less than thick. In a February 1960 letter to the journal Astronautics, Fred Singer, then science advisor to U.S. President Dwight D. Eisenhower, said of Shklovsky's theory:
If the satellite is indeed spiraling inward as deduced from astronomical observation, then there is little alternative to the hypothesis that it is hollow and therefore Martian made. The big 'if' lies in the astronomical observations; they may well be in error. Since they are based on several independent sets of measurements taken decades apart by different observers with different instruments, systematic errors may have influenced them.
Subsequently, the systematic data errors that Singer predicted were found to exist, the claim was called into doubt, and accurate measurements of the orbit available by 1969 showed that the discrepancy did not exist. Singer's critique was justified when earlier studies were discovered to have used an overestimated value of per year for the rate of altitude loss, which was later revised to per year. The secular acceleration is now attributed to tidal effects, which create drag on the moon and therefore cause it to spiral inward.
The density of Phobos has now been directly measured by spacecraft to be . Current observations are consistent with Phobos being a rubble pile. In addition, images obtained by the Viking probes in the 1970s clearly showed a natural object, not an artificial one. Nevertheless, mapping by the Mars Express probe and subsequent volume calculations do suggest the presence of voids and indicate that it is not a solid chunk of rock but a porous body. The porosity of Phobos was calculated to be 30% ± 5%, or a quarter to a third being empty.
Named geological features
Geological features on Phobos are named after astronomers who studied Phobos and people and places from Jonathan Swift's Gulliver's Travels.
Craters on Phobos
A number of craters have been named, and are listed in the following map and table.
Other named features
There is one named regio, Laputa Regio, and one named planitia, Lagado Planitia; both are named after places in Gulliver's Travels (the fictional Laputa, a flying island, and Lagado, imaginary capital of the fictional nation Balnibarbi). The only named ridge on Phobos is Kepler Dorsum, named after the astronomer Johannes Kepler.
Orbital characteristics
The orbital motion of Phobos has been intensively studied, making it "the best studied natural satellite in the Solar System" in terms of orbits completed. Its close orbit around Mars produces some unusual effects. With an altitude of , Phobos orbits Mars below the synchronous orbit radius, meaning that it moves around Mars faster than Mars itself rotates. Therefore, from the point of view of an observer on the surface of Mars, it rises in the west, moves comparatively rapidly across the sky (in 4 h 15 min or less) and sets in the east, approximately twice each Martian day (every 11 h 6 min). Because it is close to the surface and in an equatorial orbit, it cannot be seen above the horizon from latitudes greater than 70.4°. Its orbit is so low that its angular diameter, as seen by an observer on Mars, varies visibly with its position in the sky. Seen at the horizon, Phobos is about 0.14° wide; at zenith, it is 0.20°, one-third as wide as the full Moon as seen from Earth. By comparison, the Sun has an apparent size of about 0.35° in the Martian sky. Phobos' phases, inasmuch as they can be observed from Mars, take 0.3191 days (Phobos' synodic period) to run their course, a mere 13 seconds longer than Phobos' sidereal period.
Solar transits
An observer situated on the Martian surface, in a position to observe Phobos, would see regular transits of Phobos across the Sun. Several of these transits have been photographed by the Mars Rover Opportunity. During the transits, Phobos casts a shadow on the surface of Mars; this event has been photographed by several spacecraft. Phobos is not large enough to cover the Sun's disk, and so cannot cause a total eclipse.
Predicted destruction
Tidal deceleration is gradually decreasing the orbital radius of Phobos by approximately every 100 years, and with decreasing orbital radius the likelihood of breakup due to tidal forces increases, estimated in approximately 30–50 million years, or about 43 million years in one study's estimate.
Phobos' grooves were long thought to be fractures caused by the impact that formed the Stickney crater. Other modelling suggested since the 1970s support the idea that the grooves are more like "stretch marks" that occur when Phobos gets deformed by tidal forces, but in 2015 when the tidal forces were calculated and used in a new model, the stresses were too weak to fracture a solid moon of that size, unless Phobos is a rubble pile surrounded by a layer of powdery regolith about thick. Stress fractures calculated for this model line up with the grooves on Phobos. The model is supported with the discovery that some of the grooves are younger than others, implying that the process that produces the grooves is ongoing.
Given Phobos' irregular shape and assuming that it is a pile of rubble (specifically a Mohr–Coulomb body), it will eventually break up due to tidal forces when it reaches approximately 2.1 Mars radii. When Phobos is broken up, it will form a planetary ring around Mars. This predicted ring may last from 1 million to 100 million years. The fraction of the mass of Phobos that will form the ring depends on the unknown internal structure of Phobos. Loose, weakly bound material will form the ring. Components of Phobos with strong cohesion will escape tidal breakup and will enter the Martian atmosphere. It is predicted that within 30 to 50 million years it will either collide with the planet or break up into a planetary ring.
Origin
The origin of the Martian moons has been disputed. Phobos and Deimos both have much in common with carbonaceous C-type asteroids, with spectra, albedo, and density very similar to those of C- or D-type asteroids. Based on their similarity, one hypothesis is that both moons may be captured main-belt asteroids. Both moons have very circular orbits which lie almost exactly in Mars' equatorial plane, and hence a capture origin requires a mechanism for circularizing the initially highly eccentric orbit, and adjusting its inclination into the equatorial plane, most probably by a combination of atmospheric drag and tidal forces, although it is not clear that sufficient time is available for this to occur for Deimos. Capture also requires dissipation of energy. The current Martian atmosphere is too thin to capture a Phobos-sized object by atmospheric braking. Geoffrey A. Landis has pointed out that the capture could have occurred if the original body was a binary asteroid that separated under tidal forces.
Phobos could be a second-generation Solar System object that coalesced in orbit after Mars formed, rather than forming concurrently out of the same birth cloud as Mars.
Another hypothesis is that Mars was once surrounded by many Phobos- and Deimos-sized bodies, perhaps ejected into orbit around it by a collision with a large planetesimal. The high porosity of the interior of Phobos (based on the density of 1.88 g/cm3, voids are estimated to comprise 25 to 35 percent of Phobos' volume) is inconsistent with an asteroidal origin. Observations of Phobos in the thermal infrared suggest a composition containing mainly phyllosilicates, which are well known from the surface of Mars. The spectra are distinct from those of all classes of chondrite meteorites, again pointing away from an asteroidal origin. Both sets of findings support an origin of Phobos from material ejected by an impact on Mars that reaccreted in Martian orbit, similar to the prevailing theory for the origin of Earth's moon.
Some areas of the surface are reddish in color, while others are bluish. The hypothesis is that gravity pull from Mars makes the reddish regolith move over the surface, exposing relatively fresh, unweathered and bluish material from the moon, while the regolith covering it over time has been weathered due to exposure of solar radiation. Because the blue rock differs from known Martian rock, it could contradict the theory that the moon is formed from leftover planetary material after the impact of a large object.
In February 2021, Amirhossein Bagheri (ETH Zurich), Amir Khan (ETH Zurich), Michael Efroimsky (US Naval Observatory) and their colleagues proposed a new hypothesis on the origin of the moons. By analyzing the seismic and orbital data from Mars InSight Mission and other missions, they proposed that the moons are born from disruption of a common parent body around 1 to 2.7 billion years ago. The common progenitor of Phobos and Deimos was most probably hit by another object and shattered to form both moons.
Exploration
Launched missions
Phobos has been photographed in close-up by several spacecraft whose primary mission has been to photograph Mars. The first was Mariner 7 in 1969, followed by Mariner 9 in 1971, Viking 1 in 1977, Phobos 2 in 1989 Mars Global Surveyor in 1998 and 2003, Mars Express in 2004, 2008, 2010 and 2019, and Mars Reconnaissance Orbiter in 2007 and 2008. On 25 August 2005, the Spirit rover, with an excess of energy due to wind blowing dust off of its solar panels, took several short-exposure photographs of the night sky from the surface of Mars, and was able to successfully photograph both Phobos and Deimos.
The Soviet Union undertook the Phobos program with two probes, both launched successfully in July 1988. Phobos 1 was accidentally shut down by an erroneous command from ground control issued in September 1988 and lost while the craft was still en route. Phobos 2 arrived at the Mars system in January 1989 and, after transmitting a small amount of data and imagery shortly before beginning its detailed examination of Phobos' surface, the probe abruptly ceased transmission due either to failure of the onboard computer or of the radio transmitter, already operating on backup power. Other Mars missions collected more data, but no dedicated sample return mission has been successfully performed.
The Russian Space Agency launched a sample return mission to Phobos in November 2011, called Fobos-Grunt. The return capsule also included a life science experiment of The Planetary Society, called Living Interplanetary Flight Experiment, or LIFE. A second contributor to this mission was the China National Space Administration, which supplied a surveying satellite called "Yinghuo-1", which would have been released in the orbit of Mars, and a soil-grinding and sieving system for the scientific payload of the Phobos lander. However, after achieving Earth orbit, the Fobos-Grunt probe failed to initiate subsequent burns that would have sent it to Mars. Attempts to recover the probe were unsuccessful and it crashed back to Earth in January 2012.
On 1 July 2020, the Mars orbiter of the Indian Space Research Organisation was able to capture photos of the body from 4,200 km away.
Missions considered
In 1997 and 1998, the Aladdin mission was selected as a finalist in the NASA Discovery Program. The plan was to visit both Phobos and Deimos, and launch projectiles at the satellites. The probe would collect the ejecta as it performed a slow flyby (~1 km/s). These samples would be returned to Earth for study three years later. The Principal Investigator was Dr. Carle Pieters of Brown University. The total mission cost, including launch vehicle and operations was $247.7 million. Ultimately, the mission chosen to fly was MESSENGER, a probe to Mercury.
In 2007, the European aerospace subsidiary EADS Astrium was reported to have been developing a mission to Phobos as a technology demonstrator. Astrium was involved in developing a European Space Agency plan for a sample return mission to Mars, as part of the ESA's Aurora programme, and sending a mission to Phobos with its low gravity was seen as a good opportunity for testing and proving the technologies required for an eventual sample return mission to Mars. The mission was envisioned to start in 2016, was to last for three years. The company planned to use a "mothership", which would be propelled by an ion engine, releasing a lander to the surface of Phobos. The lander would perform some tests and experiments, gather samples in a capsule, then return to the mothership and head back to Earth where the samples would be jettisoned for recovery on the surface.
Proposed missions
In 2007, the Canadian Space Agency funded a study by Optech and the Mars Institute for an uncrewed mission to Phobos known as Phobos Reconnaissance and International Mars Exploration (PRIME). A proposed landing site for the PRIME spacecraft is at the "Phobos monolith", a prominent object near Stickney crater. The PRIME mission would be composed of an orbiter and lander, and each would carry 4 instruments designed to study various aspects of Phobos' geology.
In 2008, NASA Glenn Research Center began studying a Phobos and Deimos sample return mission that would use solar electric propulsion. The study gave rise to the "Hall" mission concept, a New Frontiers-class mission under further study as of 2010.
Another concept of a sample return mission from Phobos and Deimos is OSIRIS-REx II, which would use heritage technology from the first OSIRIS-REx mission.
As of January 2013, a new Phobos Surveyor mission is currently under development by a collaboration of Stanford University, NASA's Jet Propulsion Laboratory, and the Massachusetts Institute of Technology. The mission is currently in the testing phases, and the team at Stanford plans to launch the mission between 2023 and 2033.
In March 2014, a Discovery class mission was proposed to place an orbiter in Mars orbit by 2021 to study Phobos and Deimos through a series of close flybys. The mission is called Phobos And Deimos & Mars Environment (PADME). Two other Phobos missions that were proposed for the Discovery 13 selection included a mission called Merlin, which would flyby Deimos but actually orbit and land on Phobos, and another one is Pandora which would orbit both Deimos and Phobos.
The Japanese Aerospace Exploration Agency (JAXA) unveiled on 9 June 2015 the Martian Moons Exploration (MMX), a sample return mission targeting Phobos. MMX will land and collect samples from Phobos multiple times, along with conducting Deimos flyby observations and monitoring Mars' climate. By using a corer sampling mechanism, the spacecraft aims to retrieve a minimum 10 g amount of samples. NASA, ESA, DLR, and CNES are also participating in the project, and will provide scientific instruments. The U.S. will contribute the Neutron and Gamma-Ray
Spectrometer (NGRS), and France the Near IR Spectrometer (NIRS4/MacrOmega). Although the mission has been selected for implementation and is now beyond proposal stage, formal project approval by JAXA has been postponed following the Hitomi mishap. Development and testing of key components, including the sampler, is currently ongoing. , MMX is scheduled to be launched in 2026, and will return to Earth five years later.
Russia plans to repeat Fobos-Grunt mission in the late 2020s, and the European Space Agency is assessing a sample-return mission for 2024 called Phootprint.
Human missions
Phobos has been proposed as an early target for a human mission to Mars. The teleoperation of robotic scouts on Mars by humans on Phobos could be conducted without significant time delay, and planetary protection concerns in early Mars exploration might be addressed by such an approach.
A landing on Phobos would be considerably less difficult and expensive than a landing on the surface of Mars itself. A lander bound for Mars would need to be capable of atmospheric entry and subsequent return to orbit without any support facilities, or would require the creation of support facilities in-situ. A lander instead bound for Phobos could be based on equipment designed for lunar and asteroid landings. Furthermore, due to Phobos' very weak gravity, the delta-v required to land on Phobos and return is only 80% of that required for a trip to and from the surface of the Moon.
It has been proposed that the sands of Phobos could serve as a valuable material for aerobraking during a Mars landing. A relatively small amount of chemical fuel brought from Earth could be used to lift a large amount of sand from the surface of Phobos to a transfer orbit. This sand could be released in front of a spacecraft during the descent maneuver causing a densification of the atmosphere just in front of the spacecraft.
While human exploration of Phobos could serve as a catalyst for the human exploration of Mars, it could be scientifically valuable in its own right.
Space elevator base
First discussed in fiction in 1956 by Fontenay, Phobos has been proposed as a future site for space elevator construction. This would involve a pair of space elevators: one extending 6,000 km from the Mars-facing side to the edge of Mars' atmosphere, the other extending from the other side and away from Mars. A spacecraft launching from Mars' surface to the lower space elevator would only need a delta-v of , as opposed to the over needed to launch to low Mars orbit. The spacecraft could be lifted up using electrical power and then released from the upper space elevator with a hyperbolic velocity of , enough to reach Earth and a significant fraction of the velocity needed to reach the asteroid belt. The space elevators could also work in reverse to help spacecraft enter the Martian system. The great mass of Phobos means that any forces from space elevator operation would have minimal effect on its orbit. Additionally, materials from Phobos could be used for space industry.
| Physical sciences | Solar System | Astronomy |
53031 | https://en.wikipedia.org/wiki/Stefan%E2%80%93Boltzmann%20law | Stefan–Boltzmann law | The Stefan–Boltzmann law, also known as Stefan's law, describes the intensity of the thermal radiation emitted by matter in terms of that matter's temperature. It is named for Josef Stefan, who empirically derived the relationship, and Ludwig Boltzmann who derived the law theoretically.
For an ideal absorber/emitter or black body, the Stefan–Boltzmann law states that the total energy radiated per unit surface area per unit time (also known as the radiant exitance) is directly proportional to the fourth power of the black body's temperature, :
The constant of proportionality, , is called the Stefan–Boltzmann constant. It has the value
In the general case, the Stefan–Boltzmann law for radiant exitance takes the form:
where is the emissivity of the surface emitting the radiation. The emissivity is generally between zero and one. An emissivity of one corresponds to a black body.
Detailed explanation
The radiant exitance (previously called radiant emittance), , has dimensions of energy flux (energy per unit time per unit area), and the SI units of measure are joules per second per square metre (J⋅s⋅m), or equivalently, watts per square metre (W⋅m). The SI unit for absolute temperature, , is the kelvin (K).
To find the total power, , radiated from an object, multiply the radiant exitance by the object's surface area, :
Matter that does not absorb all incident radiation emits less total energy than a black body. Emissions are reduced by a factor , where the emissivity, , is a material property which, for most matter, satisfies . Emissivity can in general depend on wavelength, direction, and polarization. However, the emissivity which appears in the non-directional form of the Stefan–Boltzmann law is the hemispherical total emissivity, which reflects emissions as totaled over all wavelengths, directions, and polarizations.
The form of the Stefan–Boltzmann law that includes emissivity is applicable to all matter, provided that matter is in a state of local thermodynamic equilibrium (LTE) so that its temperature is well-defined. (This is a trivial conclusion, since the emissivity, , is defined to be the quantity that makes this equation valid. What is non-trivial is the proposition that , which is a consequence of Kirchhoff's law of thermal radiation.)
A so-called grey body is a body for which the spectral emissivity is independent of wavelength, so that the total emissivity, , is a constant. In the more general (and realistic) case, the spectral emissivity depends on wavelength. The total emissivity, as applicable to the Stefan–Boltzmann law, may be calculated as a weighted average of the spectral emissivity, with the blackbody emission spectrum serving as the weighting function. It follows that if the spectral emissivity depends on wavelength then the total emissivity depends on the temperature, i.e., . However, if the dependence on wavelength is small, then the dependence on temperature will be small as well.
Wavelength- and subwavelength-scale particles, metamaterials, and other nanostructures are not subject to ray-optical limits and may be designed to have an emissivity greater than 1.
In national and international standards documents, the symbol is recommended to denote radiant exitance; a superscript circle (°) indicates a term relate to a black body. (A subscript "e" is added when it is important to distinguish the energetic (radiometric) quantity radiant exitance, , from the analogous human vision (photometric) quantity, luminous exitance, denoted .) In common usage, the symbol used for radiant exitance (often called radiant emittance) varies among different texts and in different fields.
The Stefan–Boltzmann law may be expressed as a formula for radiance as a function of temperature. Radiance is measured in watts per square metre per steradian (W⋅m⋅sr). The Stefan–Boltzmann law for the radiance of a black body is:
The Stefan–Boltzmann law expressed as a formula for radiation energy density is:
where is the speed of light.
History
In 1864, John Tyndall presented measurements of the infrared emission by a platinum filament and the corresponding color of the filament.
The proportionality to the fourth power of the absolute temperature was deduced by Josef Stefan (1835–1893) in 1877 on the basis of Tyndall's experimental measurements, in the article Über die Beziehung zwischen der Wärmestrahlung und der Temperatur (On the relationship between thermal radiation and temperature) in the Bulletins from the sessions of the Vienna Academy of Sciences.
A derivation of the law from theoretical considerations was presented by Ludwig Boltzmann (1844–1906) in 1884, drawing upon the work of Adolfo Bartoli.
Bartoli in 1876 had derived the existence of radiation pressure from the principles of thermodynamics. Following Bartoli, Boltzmann considered an ideal heat engine using electromagnetic radiation instead of an ideal gas as working matter.
The law was almost immediately experimentally verified. Heinrich Weber in 1888 pointed out deviations at higher temperatures, but perfect accuracy within measurement uncertainties was confirmed up to temperatures of 1535 K by 1897.
The law, including the theoretical prediction of the Stefan–Boltzmann constant as a function of the speed of light, the Boltzmann constant and the Planck constant, is a direct consequence of Planck's law as formulated in 1900.
Stefan–Boltzmann constant
The Stefan–Boltzmann constant, , is derived from other known physical constants:
where is the Boltzmann constant, the is the Planck constant, and is the speed of light in vacuum.
As of the 2019 revision of the SI, which establishes exact fixed values for , , and , the Stefan–Boltzmann constant is exactly:
Thus,
Prior to this, the value of was calculated from the measured value of the gas constant.
The numerical value of the Stefan–Boltzmann constant is different in other systems of units, as shown in the table below.
Examples
Temperature of the Sun
With his law, Stefan also determined the temperature of the Sun's surface. He inferred from the data of Jacques-Louis Soret (1827–1890) that the energy flux density from the Sun is 29 times greater than the energy flux density of a certain warmed metal lamella (a thin plate). A round lamella was placed at such a distance from the measuring device that it would be seen at the same angular diameter as the Sun. Soret estimated the temperature of the lamella to be approximately 1900 °C to 2000 °C. Stefan surmised that 1/3 of the energy flux from the Sun is absorbed by the Earth's atmosphere, so he took for the correct Sun's energy flux a value 3/2 times greater than Soret's value, namely 29 × 3/2 = 43.5.
Precise measurements of atmospheric absorption were not made until 1888 and 1904. The temperature Stefan obtained was a median value of previous ones, 1950 °C and the absolute thermodynamic one 2200 K. As 2.574 = 43.5, it follows from the law that the temperature of the Sun is 2.57 times greater than the temperature of the lamella, so Stefan got a value of 5430 °C or 5700 K. This was the first sensible value for the temperature of the Sun. Before this, values ranging from as low as 1800 °C to as high as were claimed. The lower value of 1800 °C was determined by Claude Pouillet (1790–1868) in 1838 using the Dulong–Petit law. Pouillet also took just half the value of the Sun's correct energy flux.
Temperature of stars
The temperature of stars other than the Sun can be approximated using a similar means by treating the emitted energy as a black body radiation. So:
where is the luminosity, is the Stefan–Boltzmann constant, is the stellar radius and is the effective temperature. This formula can then be rearranged to calculate the temperature:
or alternatively the radius:
The same formulae can also be simplified to compute the parameters relative to the Sun:
where is the solar radius, and so forth. They can also be rewritten in terms of the surface area A and radiant exitance :
where and
With the Stefan–Boltzmann law, astronomers can easily infer the radii of stars. The law is also met in the thermodynamics of black holes in so-called Hawking radiation.
Effective temperature of the Earth
Similarly we can calculate the effective temperature of the Earth T⊕ by equating the energy received from the Sun and the energy radiated by the Earth, under the black-body approximation (Earth's own production of energy being small enough to be negligible). The luminosity of the Sun, L⊙, is given by:
At Earth, this energy is passing through a sphere with a radius of a0, the distance between the Earth and the Sun, and the irradiance (received power per unit area) is given by
The Earth has a radius of R⊕, and therefore has a cross-section of . The radiant flux (i.e. solar power) absorbed by the Earth is thus given by:
Because the Stefan–Boltzmann law uses a fourth power, it has a stabilizing effect on the exchange and the flux emitted by Earth tends to be equal to the flux absorbed, close to the steady state where:
T⊕ can then be found:
where T⊙ is the temperature of the Sun, R⊙ the radius of the Sun, and a0 is the distance between the Earth and the Sun. This gives an effective temperature of 6 °C on the surface of the Earth, assuming that it perfectly absorbs all emission falling on it and has no atmosphere.
The Earth has an albedo of 0.3, meaning that 30% of the solar radiation that hits the planet gets scattered back into space without absorption. The effect of albedo on temperature can be approximated by assuming that the energy absorbed is multiplied by 0.7, but that the planet still radiates as a black body (the latter by definition of effective temperature, which is what we are calculating). This approximation reduces the temperature by a factor of 0.71/4, giving .
The above temperature is Earth's as seen from space, not ground temperature but an average over all emitting bodies of Earth from surface to high altitude. Because of the greenhouse effect, the Earth's actual average surface temperature is about , which is higher than the effective temperature, and even higher than the temperature that a black body would have.
In the above discussion, we have assumed that the whole surface of the earth is at one temperature. Another interesting question is to ask what the temperature of a blackbody surface on the earth would be assuming that it reaches equilibrium with the sunlight falling on it. This of course depends on the angle of the sun on the surface and on how much air the sunlight has gone through. When the sun is at the zenith and the surface is horizontal, the irradiance can be as high as 1120 W/m2. The Stefan–Boltzmann law then gives a temperature of
or . (Above the atmosphere, the result is even higher: .) We can think of the earth's surface as "trying" to reach equilibrium temperature during the day, but being cooled by the atmosphere, and "trying" to reach equilibrium with starlight and possibly moonlight at night, but being warmed by the atmosphere.
Origination
Thermodynamic derivation of the energy density
The fact that the energy density of the box containing radiation is proportional to can be derived using thermodynamics. This derivation uses the relation between the radiation pressure p and the internal energy density , a relation that can be shown using the form of the electromagnetic stress–energy tensor. This relation is:
Now, from the fundamental thermodynamic relation
we obtain the following expression, after dividing by and fixing :
The last equality comes from the following Maxwell relation:
From the definition of energy density it follows that
where the energy density of radiation only depends on the temperature, therefore
Now, the equality is
after substitution of
Meanwhile, the pressure is the rate of momentum change per unit area. Since the momentum of a photon is the same as the energy divided by the speed of light,
where the factor 1/3 comes from the projection of the momentum transfer onto the normal to the wall of the container.
Since the partial derivative can be expressed as a relationship between only and (if one isolates it on one side of the equality), the partial derivative can be replaced by the ordinary derivative. After separating the differentials the equality becomes
which leads immediately to , with as some constant of integration.
Derivation from Planck's law
The law can be derived by considering a small flat black body surface radiating out into a half-sphere. This derivation uses spherical coordinates, with θ as the zenith angle and φ as the azimuthal angle; and the small flat blackbody surface lies on the xy-plane, where θ = /2.
The intensity of the light emitted from the blackbody surface is given by Planck's law,
where
is the amount of power per unit surface area per unit solid angle per unit frequency emitted at a frequency by a black body at temperature T.
is the Planck constant
is the speed of light, and
is the Boltzmann constant.
The quantity is the power radiated by a surface of area A through a solid angle in the frequency range between and .
The Stefan–Boltzmann law gives the power emitted per unit area of the emitting body,
Note that the cosine appears because black bodies are Lambertian (i.e. they obey Lambert's cosine law), meaning that the intensity observed along the sphere will be the actual intensity times the cosine of the zenith angle.
To derive the Stefan–Boltzmann law, we must integrate over the half-sphere and integrate from 0 to ∞.
Then we plug in for I:
To evaluate this integral, do a substitution,
which gives:
The integral on the right is standard and goes by many names: it is a particular case of a Bose–Einstein integral, the polylogarithm, or the Riemann zeta function . The value of the integral is (where is the Gamma function), giving the result that, for a perfect blackbody surface:
Finally, this proof started out only considering a small flat surface. However, any differentiable surface can be approximated by a collection of small flat surfaces. So long as the geometry of the surface does not cause the blackbody to reabsorb its own radiation, the total energy radiated is just the sum of the energies radiated by each surface; and the total surface area is just the sum of the areas of each surface—so this law holds for all convex blackbodies, too, so long as the surface has the same temperature throughout. The law extends to radiation from non-convex bodies by using the fact that the convex hull of a black body radiates as though it were itself a black body.
Energy density
The total energy density U can be similarly calculated, except the integration is over the whole sphere and there is no cosine, and the energy flux (U c) should be divided by the velocity c to give the energy density U:
Thus is replaced by , giving an extra factor of 4.
Thus, in total:
The product is sometimes known as the radiation constant or radiation density constant.
Decomposition in terms of photons
The Stephan–Boltzmann law can be expressed as
where the flux of photons, , is given by
and the average energy per photon,, is given by
Marr and Wilkin (2012) recommend that students be taught about instead of being taught Wien's displacement law, and that the above decomposition be taught when the Stefan–Boltzmann law is taught.
| Physical sciences | Thermodynamics | Physics |
53039 | https://en.wikipedia.org/wiki/Key%20%28cryptography%29 | Key (cryptography) | A key in cryptography is a piece of information, usually a string of numbers or letters that are stored in a file, which, when processed through a cryptographic algorithm, can encode or decode cryptographic data. Based on the used method, the key can be different sizes and varieties, but in all cases, the strength of the encryption relies on the security of the key being maintained. A key's security strength is dependent on its algorithm, the size of the key, the generation of the key, and the process of key exchange.
Scope
The key is what is used to encrypt data from plaintext to ciphertext. There are different methods for utilizing keys and encryption.
Symmetric cryptography
Symmetric cryptography refers to the practice of the same key being used for both encryption and decryption.
Asymmetric cryptography
Asymmetric cryptography has separate keys for encrypting and decrypting. These keys are known as the public and private keys, respectively.
Purpose
Since the key protects the confidentiality and integrity of the system, it is important to be kept secret from unauthorized parties. With public key cryptography, only the private key must be kept secret, but with symmetric cryptography, it is important to maintain the confidentiality of the key. Kerckhoff's principle states that the entire security of the cryptographic system relies on the secrecy of the key.
Key sizes
Key size is the number of bits in the key defined by the algorithm. This size defines the upper bound of the cryptographic algorithm's security. The larger the key size, the longer it will take before the key is compromised by a brute force attack. Since perfect secrecy is not feasible for key algorithms, researches are now more focused on computational security.
In the past, keys were required to be a minimum of 40 bits in length, however, as technology advanced, these keys were being broken quicker and quicker. As a response, restrictions on symmetric keys were enhanced to be greater in size.
Currently, 2048 bit RSA is commonly used, which is sufficient for current systems. However, current key sizes would all be cracked quickly with a powerful quantum computer.
“The keys used in public key cryptography have some mathematical structure. For example, public keys used in the RSA system are the product of two prime numbers. Thus public key systems require longer key lengths than symmetric systems for an equivalent level of security. 3072 bits is the suggested key length for systems based on factoring and integer discrete logarithms which aim to have security equivalent to a 128 bit symmetric cipher.”
Key generation
To prevent a key from being guessed, keys need to be generated randomly and contain sufficient entropy. The problem of how to safely generate random keys is difficult and has been addressed in many ways by various cryptographic systems. A key can directly be generated by using the output of a Random Bit Generator (RBG), a system that generates a sequence of unpredictable and unbiased bits. A RBG can be used to directly produce either a symmetric key or the random output for an asymmetric key pair generation. Alternatively, a key can also be indirectly created during a key-agreement transaction, from another key or from a password.
Some operating systems include tools for "collecting" entropy from the timing of unpredictable operations such as disk drive head movements. For the production of small amounts of keying material, ordinary dice provide a good source of high-quality randomness.
Establishment scheme
The security of a key is dependent on how a key is exchanged between parties. Establishing a secured communication channel is necessary so that outsiders cannot obtain the key. A key establishment scheme (or key exchange) is used to transfer an encryption key among entities. Key agreement and key transport are the two types of a key exchange scheme that are used to be remotely exchanged between entities . In a key agreement scheme, a secret key, which is used between the sender and the receiver to encrypt and decrypt information, is set up to be sent indirectly. All parties exchange information (the shared secret) that permits each party to derive the secret key material. In a key transport scheme, encrypted keying material that is chosen by the sender is transported to the receiver. Either symmetric key or asymmetric key techniques can be used in both schemes.
The Diffie–Hellman key exchange and Rivest-Shamir-Adleman (RSA) are the most two widely used key exchange algorithms. In 1976, Whitfield Diffie and Martin Hellman constructed the Diffie–Hellman algorithm, which was the first public key algorithm. The Diffie–Hellman key exchange protocol allows key exchange over an insecure channel by electronically generating a shared key between two parties. On the other hand, RSA is a form of the asymmetric key system which consists of three steps: key generation, encryption, and decryption.
Key confirmation delivers an assurance between the key confirmation recipient and provider that the shared keying materials are correct and established. The National Institute of Standards and Technology recommends key confirmation to be integrated into a key establishment scheme to validate its implementations.
Management
Key management concerns the generation, establishment, storage, usage and replacement of cryptographic keys. A key management system (KMS) typically includes three steps of establishing, storing and using keys. The base of security for the generation, storage, distribution, use and destruction of keys depends on successful key management protocols.
Key vs password
A password is a memorized series of characters including letters, digits, and other special symbols that are used to verify identity. It is often produced by a human user or a password management software to protect personal and sensitive information or generate cryptographic keys. Passwords are often created to be memorized by users and may contain non-random information such as dictionary words. On the other hand, a key can help strengthen password protection by implementing a cryptographic algorithm which is difficult to guess or replace the password altogether. A key is generated based on random or pseudo-random data and can often be unreadable to humans.
A password is less safe than a cryptographic key due to its low entropy, randomness, and human-readable properties. However, the password may be the only secret data that is accessible to the cryptographic algorithm for information security in some applications such as securing information in storage devices. Thus, a deterministic algorithm called a key derivation function (KDF) uses a password to generate the secure cryptographic keying material to compensate for the password's weakness. Various methods such as adding a salt or key stretching may be used in the generation.
| Technology | Computer security | null |
53042 | https://en.wikipedia.org/wiki/Symmetric-key%20algorithm | Symmetric-key algorithm | Symmetric-key algorithms are algorithms for cryptography that use the same cryptographic keys for both the encryption of plaintext and the decryption of ciphertext. The keys may be identical, or there may be a simple transformation to go between the two keys. The keys, in practice, represent a shared secret between two or more parties that can be used to maintain a private information link. The requirement that both parties have access to the secret key is one of the main drawbacks of symmetric-key encryption, in comparison to public-key encryption (also known as asymmetric-key encryption). However, symmetric-key encryption algorithms are usually better for bulk encryption. With exception of the one-time pad they have a smaller key size, which means less storage space and faster transmission. Due to this, asymmetric-key encryption is often used to exchange the secret key for symmetric-key encryption.
Types
Symmetric-key encryption can use either stream ciphers or block ciphers.
Stream ciphers encrypt the digits (typically bytes), or letters (in substitution ciphers) of a message one at a time. An example is ChaCha20. Substitution ciphers are well-known ciphers, but can be easily decrypted using a frequency table.
Block ciphers take a number of bits and encrypt them in a single unit, padding the plaintext to achieve a multiple of the block size. The Advanced Encryption Standard (AES) algorithm, approved by NIST in December 2001, uses 128-bit blocks.
Implementations
Examples of popular symmetric-key algorithms include Twofish, Serpent, AES (Rijndael), Camellia, Salsa20, ChaCha20, Blowfish, CAST5, Kuznyechik, RC4, DES, 3DES, Skipjack, Safer, and IDEA.
Use as a cryptographic primitive
Symmetric ciphers are commonly used to achieve other cryptographic primitives than just encryption.
Encrypting a message does not guarantee that it will remain unchanged while encrypted. Hence, often a message authentication code is added to a ciphertext to ensure that changes to the ciphertext will be noted by the receiver. Message authentication codes can be constructed from an AEAD cipher (e.g. AES-GCM).
However, symmetric ciphers cannot be used for non-repudiation purposes except by involving additional parties. See the ISO/IEC 13888-2 standard.
Another application is to build hash functions from block ciphers. See one-way compression function for descriptions of several such methods.
Construction of symmetric ciphers
Many modern block ciphers are based on a construction proposed by Horst Feistel. Feistel's construction makes it possible to build invertible functions from other functions that are themselves not invertible.
Security of symmetric ciphers
Symmetric ciphers have historically been susceptible to known-plaintext attacks, chosen-plaintext attacks, differential cryptanalysis and linear cryptanalysis. Careful construction of the functions for each round can greatly reduce the chances of a successful attack. It is also possible to increase the key length or the rounds in the encryption process to better protect against attack. This, however, tends to increase the processing power and decrease the speed at which the process runs due to the amount of operations the system needs to do.
Most modern symmetric-key algorithms appear to be resistant to the threat of post-quantum cryptography. Quantum computers would exponentially increase the speed at which these ciphers can be decoded; notably, Grover's algorithm would take the square-root of the time traditionally required for a brute-force attack, although these vulnerabilities can be compensated for by doubling key length. For example, a 128 bit AES cipher would not be secure against such an attack as it would reduce the time required to test all possible iterations from over 10 quintillion years to about six months. By contrast, it would still take a quantum computer the same amount of time to decode a 256 bit AES cipher as it would a conventional computer to decode a 128 bit AES cipher. For this reason, AES-256 is believed to be "quantum resistant".
Key management
Key establishment
Symmetric-key algorithms require both the sender and the recipient of a message to have the same secret key. All early cryptographic systems required either the sender or the recipient to somehow receive a copy of that secret key over a physically secure channel.
Nearly all modern cryptographic systems still use symmetric-key algorithms internally to encrypt the bulk of the messages, but they eliminate the need for a physically secure channel by using Diffie–Hellman key exchange or some other public-key protocol to securely come to agreement on a fresh new secret key for each session/conversation (forward secrecy).
Key generation
When used with asymmetric ciphers for key transfer, pseudorandom key generators are nearly always used to generate the symmetric cipher session keys. However, lack of randomness in those generators or in their initialization vectors is disastrous and has led to cryptanalytic breaks in the past. Therefore, it is essential that an implementation use a source of high entropy for its initialization.
Reciprocal cipher
A reciprocal cipher is a cipher where, just as one enters the plaintext into the cryptography system to get the ciphertext, one could enter the ciphertext into the same place in the system to get the plaintext. A reciprocal cipher is also sometimes referred as self-reciprocal cipher.
Practically all mechanical cipher machines implement a reciprocal cipher, a mathematical involution on each typed-in letter.
Instead of designing two kinds of machines, one for encrypting and one for decrypting, all the machines can be identical and can be set up (keyed) the same way.
Examples of reciprocal ciphers include:
Atbash
Beaufort cipher
Enigma machine
Marie Antoinette and Axel von Fersen communicated with a self-reciprocal cipher.
the Porta polyalphabetic cipher is self-reciprocal.
Purple cipher
RC4
ROT13
XOR cipher
Vatsyayana cipher
The majority of all modern ciphers can be classified as either a stream cipher, most of which use a reciprocal XOR cipher combiner, or a block cipher, most of which use a Feistel cipher or Lai–Massey scheme with a reciprocal transformation in each round.
| Technology | Computer security | null |
53045 | https://en.wikipedia.org/wiki/File%20Allocation%20Table | File Allocation Table | File Allocation Table (FAT) is a file system developed for personal computers and was the default filesystem for the MS-DOS and Windows 9x operating systems. Originally developed in 1977 for use on floppy disks, it was adapted for use on hard disks and other devices. The increase in disk drive capacity over time drove modifications to the design that resulted in versions: FAT12, FAT16, FAT32, and exFAT. FAT was replaced with NTFS as the default file system on Microsoft operating systems starting with Windows XP. Nevertheless, FAT continues to be commonly used on relatively small capacity solid-state storage technologies such as SD card, MultiMediaCard (MMC) and eMMC because of its compatibility and ease of implementation.
Uses
Historical
FAT was used on hard disks throughout the DOS and Windows 9x eras. Microsoft introduced NTFS with the Windows NT platform in 1993, but FAT remained the standard for the home user until the introduction of Windows XP in 2001. Windows Me was the final version of Windows to use FAT as its default file system.
For floppy disks, FAT has been standardized as ECMA-107 and ISO/IEC 9293:1994 (superseding ISO 9293:1987). These standards cover FAT12 and FAT16 with only short 8.3 filename support; long filenames with VFAT were partially patented. While FAT12 is used on floppy disks, FAT16 and FAT32 are typically found on the larger media.
Modern
FAT is used internally for the EFI system partition in the boot stage of EFI-compliant computers.
FAT is still used in drives expected to be used by multiple operating systems, such as in shared Windows, Linux and DOS environments. Microsoft Windows additionally comes with a pre-installed tool to convert a FAT file system into NTFS directly without the need to rewrite all files, though this cannot be reversed easily. The FAT file system is used in removable media such as floppy disks, super-floppies, memory and flash memory cards or USB flash drives. FAT is supported by portable devices such as PDAs, digital cameras, camcorders, media players, and mobile phones.
The DCF file system adopted by almost all digital cameras since 1998 defines a logical file system with 8.3 filenames and makes the use of either FAT12, FAT16, FAT32 or exFAT mandatory for its physical layer for compatibility.
Technical details
The file system uses an index table stored on the device to identify chains of data storage areas associated with a file, the File Allocation Table (FAT). The FAT is statically allocated at the time of formatting. The table is a linked list of entries for each cluster, a contiguous area of disk storage. Each entry contains either the number of the next cluster in the file, or else a marker indicating the end of the file, unused disk space, or special reserved areas of the disk. The root directory of the disk contains the number of the first cluster of each file in that directory. The operating system can then traverse the FAT, looking up the cluster number of each successive part of the disk file as a cluster chain until the end of the file is reached. Sub-directories are implemented as special files containing the directory entries of their respective files.
Each entry in the FAT linked list is a fixed number of bits: 12, 16 or 32. The maximum size of a file or a disk drive that can be accessed is the product of the largest number that can be stored in the entries (less a few values reserved to indicate unallocated space or the end of a list) and the size of the disk cluster. Even if only one byte of storage is needed to extend a file, an entire cluster must be allocated to it. As a result, large numbers of small files can result in clusters being allocated that may contain mostly "empty" data to meet the minimum cluster size.
Originally designed as an 8-bit file system, the maximum number of clusters must increase as disk drive capacity increases, and so the number of bits used to identify each cluster has grown. The successive major variants of the FAT format are named after the number of table element bits: 12 (FAT12), 16 (FAT16), and 32 (FAT32).
Variants
There are several variants of the FAT file system (e.g. FAT12, FAT16 and FAT32). FAT16 refers to both the original group of FAT file systems with 16-bit wide cluster entries and also to later variants. "VFAT" is an optional extension for long file names, which can work on top of any FAT file system. Volumes using VFAT long-filenames can be read also by operating systems not supporting the VFAT extension.
Original 8-bit FAT
The original FAT file system (or FAT structure, as it was called initially) was designed and implemented by Marc McDonald, based on a series of discussions between McDonald and Bill Gates.
It was introduced with 8-bit table elements (and valid data cluster numbers up to 0xBF) in a precursor to Microsoft's Standalone Disk BASIC-80 for an 8080-based successor of the NCR 7200 model VI data-entry terminal, equipped with 8-inch (200 mm) floppy disks, in 1977 or 1978.
In 1978, Standalone Disk BASIC-80 was ported to the 8086 using an emulator on a DEC PDP-10, since no real 8086 systems were available at this time.
The FAT file system was also used in Microsoft's MDOS/MIDAS, an operating system for 8080/Z80 platforms written by McDonald since 1979.
The Standalone Disk BASIC version supported three FATs, whereas this was a parameter for MIDAS. Reportedly, MIDAS was also prepared to support 10-bit, 12-bit and 16-bit FAT variants. While the size of directory entries was 16 bytes in Standalone Disk BASIC, MIDAS instead occupied 32 bytes per entry.
Tim Paterson of Seattle Computer Products (SCP) was first introduced to Microsoft's FAT structure when he helped Bob O'Rear adapting the Standalone Disk BASIC-86 emulator port onto SCP's S-100 bus 8086 CPU board prototype during a guest week at Microsoft in May 1979. The final product was shown at Lifeboat Associates' booth stand at the National Computer Conference in New York on June 4–7, 1979, where Paterson learned about the more sophisticated FAT implementation in MDOS/MIDAS and McDonald talked to him about the design of the file system.
FAT12
Between April and August 1980, while borrowing the FAT concept for SCP's own 8086 operating system QDOS 0.10, Tim Paterson extended the table elements to 12 bits, reduced the number of FATs to two, redefined the semantics of some of the reserved cluster values, and modified the disk layout, so that the root directory was now located between the FAT and the data area for his implementation of FAT12. Paterson also increased the nine-character (6.3) filename length limit to eleven characters to support CP/M-style 8.3 filenames and File Control Blocks. The format used in Microsoft Standalone Disk BASIC's 8-bit file system precursor was not supported by QDOS. By August 1980, QDOS had been renamed to 86-DOS. Starting with 86-DOS 0.42, the size and layout of directory entries was changed from 16 bytes to 32 bytes in order to add a file date stamp and increase the theoretical file size limit beyond the previous limit of 16 MB.
86-DOS 1.00 became available in early 1981. Later in 1981, 86-DOS evolved into Microsoft's MS-DOS and IBM PC DOS.
The capability to read previously formatted volumes with 16-byte directory entries was dropped with MS-DOS 1.20.
FAT12 used 12-bit entries for the cluster addresses; some values were reserved to mark the end of a chain of clusters, to mark unusable areas of the disk, or for other purposes, so the maximum number of clusters was limited to 4078. To conserve disk space, two 12-bit FAT entries used three consecutive 8-bit bytes on disk, requiring manipulation to unpack the 12-bit values. This was sufficient for the original floppy disk drives, and small hard disks up to 32 megabytes. The FAT16B version available with DOS 3.31 supported 32-bit sector numbers, and so increased the volume size limit.
All the control structures fit inside the first track, to avoid head movement during read and write operations. Any bad sector in the control structures area would make the disk unusable. The DOS formatting tool rejected such disks completely. Bad sectors were allowed only in the file data area. Clusters containing bad sectors were marked unusable with the reserved value 0xFF7 in the FAT.
While 86-DOS supported three disk formats (250.25 KB, 616 KB and 1232 KB, with FAT IDs 0xFF and 0xFE) on 8-inch (200 mm) floppy drives, IBM PC DOS 1.0, released with the original IBM Personal Computer in 1981, supported only an 8-sector floppy format with a formatted capacity of 160 KB (FAT ID 0xFE) for single-sided 5.25-inch floppy drives, and PC DOS 1.1 added support for a double-sided format with 320 KB (FAT ID 0xFF). PC DOS 2.0 introduced support for 9-sector floppy formats with 180 KB (FAT ID 0xFC) and 360 KB (FAT ID 0xFD).
86-DOS 1.00 and PC DOS 1.0 directory entries included only one date, the last modified date. PC DOS 1.1 added the last modified time. PC DOS 1.x file attributes included a hidden bit and system bit, with the remaining six bits undefined. At this time, DOS did not support sub-directories, but typically there were only a few dozen files on a diskette.
The PC XT was the first PC with an IBM-supplied hard drive, and PC DOS 2.0 supported that hard drive with FAT12 (FAT ID 0xF8). The fixed assumption of 8 sectors per clusters on hard disks practically limited the maximum partition size to 16 MB for 512 byte sectors and 4 KB clusters.
The BIOS Parameter Block (BPB) was introduced with PC DOS 2.0 as well, and this version also added read-only, archive, volume label, and directory attribute bits for hierarchical sub-directories.
MS-DOS 3.0 introduced support for high-density 1.2 MB 5.25-inch diskettes (media descriptor 0xF9), which notably had 15 sectors per track, hence more space for the FATs.
FAT12 remains in use on all common floppy disks, including 1.44 MB and later 2.88 MB disks (media descriptor byte 0xF0).
Initial FAT16
In 1984, IBM released the PC AT, which required PC DOS 3.0 to access its 20 MB hard disk. Microsoft introduced MS-DOS 3.0 in parallel. Cluster addresses were increased to 16-bit, allowing for up to 65,526 clusters per volume. However, the maximum possible number of sectors and the maximum partition size of 32 MB did not change. Although cluster addresses were 16 bits, this format was not what today is commonly understood as FAT16.
A partition type 0x04 indicates this form of FAT16 with less than 65,536 sectors (less than 32 MB for sector size 512). The benefit of FAT16 was the use of smaller clusters, making disk usage more efficient, particularly for large numbers of files only a few hundred bytes in size.
As MS-DOS 3.0 formatted all 16 MB-32 MB partitions in the FAT16 format, a 20 MB hard disk formatted under MS-DOS 3.0 was not accessible by MS-DOS 2.0. MS-DOS 3.0 to MS-DOS 3.30 could still access FAT12 partitions under 15 MB, but required all 16 MB-32 MB partitions to be FAT16, and so could not access MS-DOS 2.0 partitions in this size range. MS-DOS 3.31 and higher could access 16 MB-32 MB FAT12 partitions again.
Logical sectored FAT
MS-DOS and PC DOS implementations of FAT12 and FAT16 could not access disk partitions larger than 32 megabytes. Several manufacturers developed their own FAT variants within their OEM versions of MS-DOS.
Some vendors (AST and NEC) supported eight, instead of the standard four, primary partition entries in their custom extended Master Boot Record (MBR), and they adapted MS-DOS to use more than a single primary partition.
Other vendors worked around the volume size limits imposed by the 16-bit sector entries by increasing the apparent size of the sectors the file system operated on. These logical sectors were larger (up to 8192 bytes) than the physical sector size (still 512 bytes) on the disk. The DOS-BIOS or System BIOS would then combine multiple physical sectors into logical sectors for the file system to work with.
These changes were transparent to the file system implementation in the DOS kernel. The underlying DOS-BIOS translated these logical sectors into physical sectors according to partitioning information and the drive's physical geometry.
The drawback of this approach was increased memory used for sector buffering and deblocking. Since older DOS versions could not use large logical sectors, the OEMs introduced new partition IDs for their FAT variants in order to hide them from off-the-shelf issues of MS-DOS and PC DOS. Known partition IDs for logical sectored FATs include: 0x08 (Commodore MS-DOS 3.x), 0x11 (Leading Edge MS-DOS 3.x), 0x14 (AST MS-DOS 3.x), 0x24 (NEC MS-DOS 3.30), 0x56 (AT&T MS-DOS 3.x), 0xE5 (Tandy MS-DOS), 0xF2 (Sperry IT MS-DOS 3.x, Unisys MS-DOS 3.3 – also used by Digital Research DOS Plus 2.1). OEM versions like Toshiba MS-DOS, Wyse MS-DOS 3.2 and 3.3, as well as Zenith MS-DOS are also known to have utilized logical sectoring.
While non-standard and sub-optimal, these FAT variants are perfectly valid according to the specifications of the file system itself. Therefore, even if default issues of MS-DOS and PC DOS were not able to cope with them, most of these vendor-specific FAT12 and FAT16 variants can be mounted by more flexible file system implementations in operating systems such as DR-DOS, simply by changing the partition ID to one of the recognized types. Also, if they no longer need to be recognized by their original operating systems, existing partitions can be "converted" into FAT12 and FAT16 volumes more compliant with versions of MS-DOS/PC DOS 4.0–6.3, which do not support sector sizes different from 512 bytes, by switching to a BPB with 32-bit entry for the number of sectors, as introduced since DOS 3.31 (see FAT16B below), keeping the cluster size and reducing the logical sector size in the BPB down to 512 bytes, while at the same time increasing the counts of logical sectors per cluster, reserved logical sectors, total logical sectors, and logical sectors per FAT by the same factor.
A parallel development in MS-DOS / PC DOS which allowed an increase in the maximum possible FAT size was the introduction of multiple FAT partitions on a hard disk. To allow the use of more FAT partitions in a compatible way, a new partition type was introduced in PC DOS 3.2 (1986), the extended partition (EBR), which is a container for an additional partition called logical drive. Since PC DOS 3.3 (April 1987), there is another, optional extended partition containing the next logical drive, and so on. The MBR of a hard disk can either define up to four primary partitions, or an extended partition in addition to up to three primary partitions.
Final FAT16
In November 1987, Compaq Personal Computer DOS 3.31 (a modified OEM version of MS-DOS 3.3 released by Compaq with their machines) introduced what today is simply known as the FAT16 format, with the expansion of the 16-bit disk sector count to 32 bits in the BPB.
Although the on-disk changes were minor, the entire DOS disk driver had to be converted to use 32-bit sector numbers, a task complicated by the fact that it was written in 16-bit assembly language.
The result was initially called the DOS 3.31 Large File System. Microsoft's DSKPROBE tool refers to type 0x06 as BigFAT, whereas some older versions of FDISK described it as BIGDOS. Technically, it is known as FAT16B.
Since older versions of DOS were not designed to cope with more than 65,535 sectors, it was necessary to introduce a new partition type for this format in order to hide it from pre-3.31 issues of DOS. The original form of FAT16 (with less than 65,536 sectors) had a partition type 0x04. To deal with disks larger than this, type 0x06 was introduced to indicate 65,536 or more sectors. In addition to this, the disk driver was expanded to cope with more than 65,535 sectors as well. The only other difference between the original FAT16 and the newer FAT16B format is the usage of a newer BPB format with 32-bit sector entry. Therefore, newer operating systems supporting the FAT16B format can cope also with the original FAT16 format without any necessary changes.
If partitions to be used by pre-DOS 3.31 issues of DOS need to be created by modern tools, the only criteria theoretically necessary to meet are a sector count of less than 65536, and the usage of the old partition ID (0x04). In practice however, type 0x01 and 0x04 primary partitions should not be physically located outside the first 32 MB of the disk, due to other restrictions in MS-DOS 2.x, which could not cope with them otherwise.
In 1988, the FAT16B improvement became more generally available through DR DOS 3.31, PC DOS 4.0, OS/2 1.1, and MS-DOS 4.0. The limit on partition size was dictated by the 8-bit signed count of sectors per cluster, which originally had a maximum power-of-two value of 64. With the standard hard disk sector size of 512 bytes, this gives a maximum of 32 KB cluster size, thereby fixing the "definitive" limit for the FAT16 partition size at 2 GB for sector size 512. On magneto-optical media, which can have 1 or 2 KB sectors instead of 0.5 KB, this size limit is proportionally larger.
Much later, Windows NT increased the maximum cluster size to 64 KB, by considering the sectors-per-cluster count as unsigned. However, the resulting format was not compatible with any other FAT implementation of the time, and it generated greater internal fragmentation. Windows 98, SE and ME also supported reading and writing this variant, but its disk utilities did not work with it and some FCB services are not available for such volumes. This contributes to a confusing compatibility situation.
Prior to 1995, versions of DOS accessed the disk via CHS addressing only. When Windows 95(MS-DOS 7.0) introduced LBA disk access, partitions could start being physically located outside the first c. 8 GB of this disk and thereby out of the reach of the traditional CHS addressing scheme. Partitions partially or fully located beyond the CHS barrier therefore had to be hidden from non-LBA-enabled operating systems by using the new partition type 0x0E in the partition table instead. FAT16 partitions using this partition type are also named FAT16X. The only difference, compared to previous FAT16 partitions, is the fact that some CHS-related geometry entries in the BPB record, namely the number of sectors per track and the number of heads, may contain no or misleading values and should not be used.
The number of root directory entries available for FAT12 and FAT16 is determined when the volume is formatted, and is stored in a 16-bit field. For a given number RDE and sector size SS, the number RDS of root directory sectors is RDS = ceil((RDE × 32) / SS), and RDE is normally chosen to fill these sectors, i.e., RDE × 32 = RDS × SS. FAT12 and FAT16 media typically use 512 root directory entries on non-floppy media. Some third-party tools, like mkdosfs, allow the user to set this parameter.
FAT32
In order to overcome the volume size limit of FAT16, while at the same time allowing DOS real-mode code to handle the format, Microsoft designed a new version of the file system, FAT32, which supported an increased number of possible clusters, but could reuse most of the existing code, so that the conventional memory footprint was increased by less than 5 KB under DOS. Cluster values are represented by 32-bit numbers, of which 28 bits are used to hold the cluster number.
Maximal sizes
The FAT32 boot sector uses a 32-bit field for the sector count, limiting the maximal FAT32 volume size to 2 terabytes with a sector size of 512 bytes. The maximum FAT32 volume size is 16 TB with a sector size of 4,096 bytes. The built-in Windows shell disk format tool on Windows NT arbitrarily only supports volume sizes up to 32 GB, but Windows supports reading and writing to preexisting larger FAT32 volumes, and these can be created with the command prompt, PowerShell or third-party tools, or by formatting the volume on a non-Windows system or on a Windows 9x system with FAT32 support and then transferring it to the Windows NT system. In August 2024, Microsoft released an update to Windows 11 preview builds that allows for the creation of FAT32 partitions up to 2TB in size.
The maximal possible size for a file on a FAT32 volume is 4 GB minus 1 byte, or 4,294,967,295 (232 − 1) bytes. This limit is a consequence of the 4-byte file length entry in the directory table and would also affect relatively huge FAT16 partitions enabled by a sufficient sector size.
Like FAT12 and FAT16, FAT32 does not include direct built-in support for long filenames, but FAT32 volumes can optionally hold VFAT long filenames in addition to short filenames in exactly the same way as VFAT long filenames have been optionally implemented for FAT12 and FAT16 volumes.
Development
FAT32 was introduced with Windows 95 OSR2(MS-DOS 7.1) in 1996, although reformatting was needed to use it, and DriveSpace 3 (the version that came with Windows 95 OSR2 and Windows 98) never supported it. Windows 98 introduced a utility to convert existing hard disks from FAT16 to FAT32 without loss of data.
In the Windows NT line, native support for FAT32 arrived in Windows 2000. A free FAT32 driver for Windows NT 4.0 was available from Winternals, a company later acquired by Microsoft. The acquisition of the driver from official sources is no longer possible. Since 1998, Caldera's dynamically loadable DRFAT32 driver could be used to enable FAT32 support in DR-DOS. The first version of DR-DOS to natively support FAT32 and LBA access was OEM DR-DOS 7.04 in 1999. That same year IMS introduced native FAT32 support with REAL/32 7.90, and IBM 4690 OS added FAT32 support with version 2. Ahead Software provided another dynamically loadable FAT32.EXE driver for DR-DOS 7.03 with Nero Burning ROM in 2004. IBM introduced native FAT32 support with OEM PC DOS 7.1 in 1999.
Two partition types have been reserved for FAT32 partitions, 0x0B and 0x0C. The latter type is also named FAT32X in order to indicate usage of LBA disk access instead of CHS. On such partitions, CHS-related geometry entries, namely the CHS sector addresses in the MBR as well as the number of sectors per track and the number of heads in the EBPB record, may contain no or misleading values and should not be used.
Extensions
Extended attributes
OS/2 heavily depends on extended attributes (EAs) and stores them in a hidden file called "EA␠DATA.␠SF" in the root directory of the FAT12 or FAT16 volume. This file is indexed by two previously reserved bytes in the file's (or directory's) directory entry at offset 0x14. In the FAT32 format, these bytes hold the upper 16 bits of the starting cluster number of the file or directory, hence making it impossible to store OS/2 EAs on FAT32 using this method.
However, the third-party FAT32 installable file system (IFS) driver FAT32.IFS version 0.70 and higher by Henk Kelder & Netlabs for OS/2, eComStation and ArcaOS stores extended attributes in extra files with filenames having the string "␠EA.␠SF" appended to the regular filename of the file to which they belong. The driver also utilizes the byte at offset 0x0C in directory entries to store a special mark byte indicating the presence of extended attributes to help speed up things. (This extension is critically incompatible with the FAT32+ method to store files larger than 4 GB minus 1 on FAT32 volumes.)
Extended attributes are accessible via the Workplace Shell desktop, through REXX scripts, and many system GUI and command-line utilities (such as 4OS2).
To accommodate its OS/2 subsystem, Windows NT supports the handling of extended attributes in HPFS, NTFS, FAT12 and FAT16. It stores EAs on FAT12, FAT16 and HPFS using exactly the same scheme as OS/2, but does not support any other kind of ADS as held on NTFS volumes. Trying to copy a file with any ADS other than EAs from an NTFS volume to a FAT or HPFS volume gives a warning message with the names of the ADSs that will be lost. It does not support the FAT32.IFS method to store EAs on FAT32 volumes.
Windows 2000 onward acts exactly as Windows NT, except that it ignores EAs when copying to FAT32 without any warning (but shows the warning for other ADSs, like "Macintosh Finder Info" and "Macintosh Resource Fork").
Cygwin uses "EA␠DATA.␠SF" files as well.
Long file names
One of the user experience goals for the designers of Windows 95 was the ability to use long filenames (LFNs—up to 255 UTF-16 code units long), in addition to classic 8.3 filenames (SFNs). For backward and forward compatibility, LFNs were implemented as an optional extension on top of the existing FAT file system structures using a workaround in the way directory entries are laid out.
This transparent method to store long file names in the existing FAT file systems without altering their data structures is usually known as VFAT (for "Virtual FAT") after the Windows 95 virtual device driver.
Non VFAT-enabled operating systems can still access the files under their short file name alias without restrictions; however, the associated long file names may be lost when files with long filenames are copied under non VFAT-aware operating systems.
In Windows NT, support for VFAT long filenames began with version 3.5.
Linux provides a VFAT filesystem driver to work with FAT volumes with VFAT long filenames. For some time, a UVFAT driver was available to provide combined support for UMSDOS-style permissions with VFAT long filenames.
OS/2 added long filename support to FAT using extended attributes (EA) before the introduction of VFAT. Thus, VFAT long filenames are invisible to OS/2, and EA long filenames are invisible to Windows; therefore, experienced users of both operating systems would have to manually rename the files.
Human68K supported up to 18.3 filenames and (Shift JIS) Kanji characters in a proprietary FAT file system variant.
In order to support Java applications, the FlexOS-based IBM 4690 OS version 2 introduced its own virtual file system (VFS) architecture to store long filenames in the FAT file system in a backwards-compatible fashion. If enabled, the virtual filenames (VFN) are available under separate logical drive letters, whereas the real filenames (RFN) remain available under the original drive letters.
Forks and alternate data streams
The FAT file system itself is not designed for supporting alternate data streams (ADS), but some operating systems that heavily depend on them have devised various methods for handling them on FAT volumes. Such methods either store the additional information in extra files and directories (classic Mac OS and macOS), or give new semantics to previously unused fields of the FAT on-disk data structures (OS/2 and Windows NT).
Mac OS using PC Exchange stores its various dates, file attributes and long filenames in a hidden file called "FINDER.DAT", and resource forks (a common Mac OS ADS) in a subdirectory called "RESOURCE.FRK", in every directory where they are used. From PC Exchange 2.1 onwards, they store the Mac OS long filenames as standard FAT long filenames and convert FAT filenames longer than 31 characters to unique 31-character filenames, which can then be made visible to Macintosh applications.
macOS stores resource forks and metadata (file attributes, other ADS) using AppleDouble format in a hidden file with a name constructed from the owner filename prefixed with "._", and Finder stores some folder and file metadata in a hidden file called ".DS_Store" (but note that Finder uses .DS_Store even on macOS' native filesystem, HFS+).
UMSDOS permissions and filenames
Early Linux distributions also supported a format known as UMSDOS, a FAT variant with Unix file attributes (such as long file name and access permissions) stored in a separate file called "--linux-.---". UMSDOS fell into disuse after VFAT was released and it is not enabled by default in Linux from version 2.5.7 onwards. For some time, Linux also provided combined support for UMSDOS-style permissions and VFAT long filenames through UVFAT.
FAT+
In 2007 the open FAT+ draft proposed how to store larger files up to 256 GB minus 1 byte, or 274,877,906,943 (238 − 1) bytes, on slightly modified and otherwise backward-compatible FAT32 volumes, but imposes a risk that disk tools or FAT32 implementations not aware of this extension may truncate or delete files exceeding the normal FAT32 file size limit. Support for FAT32+ and FAT16+ is limited to some versions of DR-DOS and not available in mainstream operating systems. (This extension is critically incompatible with the /EAS option of the FAT32.IFS method to store OS/2 extended attributes on FAT32 volumes.)
Derivatives
Turbo FAT
In its NetWare File System (NWFS) Novell implemented a heavily modified variant of a FAT file system for the NetWare operating system. For larger files it utilized a performance feature named Turbo FAT.
FATX
FATX is a family of file systems designed for Microsoft's Xbox video game console hard disk drives and memory cards, introduced in 2001.
While resembling the same basic design ideas as FAT16 and FAT32, the FATX16 and FATX32 on-disk structures are simplified, but fundamentally incompatible with normal FAT16 and FAT32 file systems, making it impossible for normal FAT file system drivers to mount such volumes.
The non-bootable superblock sector is 4 KB in size and holds an 18 byte large BPB-like structure completely different from normal BPBs. Clusters are typically 16 KB in size and there is only one copy of the FAT on the Xbox. Directory entries are 64 bytes in size instead of the normal 32 bytes. Files can have filenames up to 42 characters long using the OEM character set and be up to 4 GB minus 1 byte in size. The on-disk timestamps hold creation, modification and access dates and times but differ from FAT: in FAT, the epoch is 1980; in FATX, the epoch is 2000. On the Xbox 360, the epoch is 1980.
exFAT
exFAT is a file system introduced with Windows Embedded CE 6.0 in November 2006 and brought to the Windows NT family with Vista Service Pack 1 and Windows XP Service Pack 3 (or separate installation of Windows XP Update KB955704). It is loosely based on the File Allocation Table architecture, but incompatible, proprietary and protected by patents.
exFAT is intended for use on flash drives and memory cards such as SDXC and Memory Stick XC, where FAT32 is otherwise used. Vendors usually pre-format SDXC cards with it. Its main benefit is its exceeding of the 4 GB file size limit, as file size references are stored with eight instead of four bytes, increasing the limit to 264 − 1 bytes.
Microsoft's GUI and command-line format utilities offer it as an alternative to NTFS (and, for smaller partitions, to FAT16B and FAT32). The MBR partition type is 0x07 (the same as used for IFS, HPFS, and NTFS). Logical geometry information located in the VBR is stored in a format not resembling any kind of BPB.
In early 2010, the file system was reverse-engineered by the SANS Institute. On August 28, 2019, Microsoft published the technical specification for exFAT so that it can be used in the Linux kernel and other operating systems.
Patents
Microsoft applied for, and was granted, a series of patents for key parts of the FAT file system in the mid-1990s. All four pertain to long-filename extensions to FAT first seen in Windows 95: U.S. patent 5,579,517, U.S. patent 5,745,902, U.S. patent 5,758,352, U.S. patent 6,286,013 (all expired since 2013).
On December 3, 2003, Microsoft announced that it would be offering licenses for use of its FAT specification and "associated intellectual property", at the cost of a royalty per unit sold, with a maximum royalty per license agreement. To this end, Microsoft cited four patents on the FAT file system as the basis of its intellectual property claims.
In the EFI FAT32 specification, Microsoft specifically grants a number of rights, which many readers have interpreted as permitting operating system vendors to implement FAT. Non-Microsoft patents affecting FAT include: U.S. patent 5,367,671, specific to the OS/2 extended object attributes (expired in 2011).
Challenges and lawsuits
The Public Patent Foundation (PUBPAT) submitted evidence to the US Patent and Trademark Office (USPTO) in 2004 disputing the validity of U.S. patent 5,579,517, including prior art references from Xerox and IBM. The USPTO opened an investigation and concluded by rejecting all claims in the patent. The next year, the USPTO further announced that following the re-examination process, it affirmed the rejection of '517 and additionally found U.S. patent 5,758,352 invalid on the grounds that the patent had incorrect assignees.
However, in 2006, the USPTO ruled that features of Microsoft's implementation of the FAT system were "novel and non-obvious", reversing both earlier decisions and leaving the patents valid.
In February 2009, Microsoft filed a patent infringement lawsuit against TomTom alleging that the device maker's products infringe on patents related to VFAT long filenames. As some TomTom products are based on Linux, this marked the first time that Microsoft tried to enforce its patents against the Linux platform. The lawsuit was settled out of court the following month with an agreement that Microsoft be given access to four of TomTom's patents, that TomTom will drop support for the VFAT long filenames from its products, and that in return Microsoft not seek legal action against TomTom for the five-year duration of the settlement agreement.
In October 2010, Microsoft filed a patent infringement lawsuit against Motorola alleging several patents (including two of the VFAT patents) were not licensed for use in the Android operating system. They also submitted a complaint to the ITC.
Developers of open source software have designed methods intended to circumvent Microsoft's patents.
In 2013, patent EP0618540 "common name space for long and short filenames" (expired since 2014) was invalidated in Germany. After the appeal was withdrawn, this judgment became final on the 28th October 2015.
| Technology | Data storage and memory | null |
53136 | https://en.wikipedia.org/wiki/Weapon%20of%20mass%20destruction | Weapon of mass destruction | A weapon of mass destruction (WMD) is a biological, chemical, radiological, nuclear, or any other weapon that can kill or significantly harm many people or cause great damage to artificial structures (e.g., buildings), natural structures (e.g., mountains), or the biosphere. The scope and usage of the term has evolved and been disputed, often signifying more politically than technically. Originally coined in reference to aerial bombing with chemical explosives during World War II, it has later come to refer to large-scale weaponry of warfare-related technologies, such as biological, chemical, radiological, or nuclear warfare.
Early uses of this term
The first use of the term "weapon of mass destruction" on record is by Cosmo Gordon Lang, Archbishop of Canterbury, in 1937 in reference to the bombing of Guernica, Spain:
At the time, nuclear weapons had not been developed fully. Japan conducted research on biological weapons , and chemical weapons had seen wide battlefield use in World War I. Their use was outlawed by the Geneva Protocol of 1925. Italy used mustard agent against civilians and soldiers in Ethiopia in 1935–36.
Following the atomic bombings of Hiroshima and Nagasaki that ended World War II and during the Cold War, the term came to refer more to non-conventional weapons. The application of the term to specifically nuclear and radiological weapons is traced by William Safire to the Russian phrase "Оружие массового поражения" – oruzhiye massovogo porazheniya (weapon of mass destruction).
William Safire credits James Goodby (of the Brookings Institution) with tracing what he considers the earliest known English-language use soon after the nuclear bombing of Hiroshima and Nagasaki (although it is not quite verbatim): a communique from a 15 November 1945, meeting of Harry Truman, Clement Attlee and Mackenzie King (probably drafted by Vannevar Bush, as Bush claimed in 1970) referred to "weapons adaptable to mass destruction."
Safire says Bernard Baruch used that exact phrase in 1946 (in a speech at the United Nations probably written by Herbert Bayard Swope). The phrase found its way into the very first resolution the United Nations General assembly adopted in January 1946 in London, which used the wording "the elimination from national armaments of atomic weapons and of all other weapons adaptable to mass destruction." The resolution also created the Atomic Energy Commission (predecessor of the International Atomic Energy Agency (IAEA)).
An exact use of this term was given in a lecture titled "Atomic Energy as a Contemporary Problem" by J. Robert Oppenheimer. He delivered the lecture to the Foreign Service and the State Department, on 17 September 1947.
It is a very far reaching control which would eliminate the rivalry between nations in this field, which would prevent the surreptitious arming of one nation against another, which would provide some cushion of time before atomic attack, and presumably therefore before any attack with weapons of mass destruction, and which would go a long way toward removing atomic energy at least as a source of conflict between the powers.
The term was also used in the introduction to the hugely influential U.S. government document known as NSC 68 written in 1950.
During a speech at Rice University on 12 September 1962, President John F. Kennedy spoke of not filling space "with weapons of mass destruction, but with instruments of knowledge and understanding." The following month, during a televised presentation about the Cuban Missile Crisis on 22 October 1962, Kennedy made reference to "offensive weapons of sudden mass destruction."
An early use of the exact phrase in an international treaty is in the Outer Space Treaty of 1967, but the treaty provides no definition of the phrase, and the treaty also categorically prohibits the stationing of "weapons" and the testing of "any type of weapon" in outer space, in addition to its specific prohibition against placing in orbit, or installing on celestial bodies, "any objects carrying nuclear weapons or any other kinds of weapons of mass destruction."
Evolution of its use
During the Cold War, the term "weapons of mass destruction" was primarily a reference to nuclear weapons. At the time, in the West the euphemism "strategic weapons" was used to refer to the American nuclear arsenal. However, there is no precise definition of the "strategic" category, neither considering range nor yield of the nuclear weapon.
Subsequent to Operation Opera, the destruction of a pre-operational nuclear reactor inside Iraq by the Israeli Air Force in 1981, the Israeli prime minister, Menachem Begin, countered criticism by saying that "on no account shall we permit an enemy to develop weapons of mass destruction against the people of Israel." This policy of pre-emptive action against real or perceived weapons of mass destruction became known as the Begin Doctrine.
The term "weapons of mass destruction" continued to see periodic use, usually in the context of nuclear arms control; Ronald Reagan used it during the 1986 Reykjavík Summit, when referring to the 1967 Outer Space Treaty. Reagan's successor, George H. W. Bush, used the term in a 1989 speech to the United Nations, primarily in reference to chemical arms.
The end of the Cold War reduced U.S. reliance on nuclear weapons as a deterrent, causing it to shift its focus to disarmament. With the 1990 invasion of Kuwait and 1991 Gulf War, Iraq's nuclear, biological, and chemical weapons programs became a particular concern of the first Bush Administration. Following the war, Bill Clinton and other western politicians and media continued to use the term, usually in reference to ongoing attempts to dismantle Iraq's weapons programs.
After the 11 September 2001 attacks and the 2001 anthrax attacks in the United States, an increased fear of nonconventional weapons and asymmetric warfare took hold in many countries. The fear reached a crescendo with the 2002 Iraq disarmament crisis and the alleged existence of weapons of mass destruction in Iraq that became the primary justification for the 2003 invasion of Iraq; however, American forces found none in Iraq. They found old stockpiles of chemical munitions including sarin and mustard agents, but all were considered to be unusable because of corrosion or degradation. Iraq, however, declared a chemical weapons stockpile in 2009 which U.N. personnel had secured after the 1991 Gulf War. The stockpile contained mainly chemical precursors, but some munitions remained usable.
Because of its prolific use and (worldwide) public profile during this period, the American Dialect Society voted "weapons of mass destruction" (and its abbreviation, "WMD") the word of the year in 2002, and in 2003 Lake Superior State University added WMD to its list of terms banished for "Mis-use, Over-use and General Uselessness" (and "as a card that trumps all forms of aggression").
In its criminal complaint against the main suspect of the Boston Marathon bombing of 15 April 2013, the FBI refers to a pressure-cooker improvised bomb as a "weapon of mass destruction."
There have been calls to classify at least some classes of cyber weapons as WMD, in particular those aimed to bring about large-scale (physical) destruction, such as by targeting critical infrastructure. However, some scholars have objected to classifying cyber weapons as WMD on the grounds that they "cannot [currently] directly injure or kill human beings as efficiently as guns or bombs" or clearly "meet the legal and historical definitions" of WMD.
Definitions of the term
United States
Strategic definition
The most widely used definition of "weapons of mass destruction" is that of nuclear, biological, or chemical weapons (NBC) although there is no treaty or customary international law that contains an authoritative definition. Instead, international law has been used with respect to the specific categories of weapons within WMD, and not to WMD as a whole. While nuclear, chemical and biological weapons are regarded as the three major types of WMDs, some analysts have argued that radiological materials as well as missile technology and delivery systems such as aircraft and ballistic missiles could be labeled as WMDs as well.
However, there is an argument that nuclear and biological weapons do not belong in the same category as chemical and "dirty bomb" radiological weapons, which have limited destructive potential (and close to none, as far as property is concerned), whereas nuclear and biological weapons have the unique ability to kill large numbers of people with very small amounts of material, and thus could be said to belong in a class by themselves.
The NBC definition has also been used in official U.S. documents, by the U.S. President, the U.S. Central Intelligence Agency, the U.S. Department of Defense, and the U.S. Government Accountability Office.
Other documents expand the definition of WMD to also include radiological or conventional weapons. The U.S. military refers to WMD as:
This may also refer to nuclear ICBMs (intercontinental ballistic missiles).
The significance of the words separable and divisible part of the weapon is that missiles such as the Pershing II and the SCUD are considered weapons of mass destruction, while aircraft capable of carrying bombloads are not.
In 2004, the United Kingdom's Butler Review recognized the "considerable and long-standing academic debate about the proper interpretation of the phrase 'weapons of mass destruction. The committee set out to avoid the general term but when using it, employed the definition of United Nations Security Council Resolution 687, which defined the systems which Iraq was required to abandon:
"Nuclear weapons or nuclear-weapons-usable material or any sub-systems or components or any research, development, support or manufacturing facilities relating to [nuclear weapons].
Chemical and biological weapons and all stocks of agents and all related subsystems and components and all research, development, support and manufacturing facilities.
Ballistic missiles with a range greater than 150 kilometres and related major parts, and repair and production facilities."
Chemical weapons expert Gert G. Harigel considers only nuclear weapons true weapons of mass destruction, because "only nuclear weapons are completely indiscriminate by their explosive power, heat radiation and radioactivity, and only they should therefore be called a weapon of mass destruction". He prefers to call chemical and biological weapons "weapons of terror" when aimed against civilians and "weapons of intimidation" for soldiers.
Testimony of one such soldier expresses the same viewpoint. For a period of several months in the winter of 2002–2003, U.S. Deputy Secretary of Defense Paul Wolfowitz frequently used the term "weapons of mass terror", apparently also recognizing the distinction between the psychological and the physical effects of many things currently falling into the WMD category.
Gustavo Bell Lemus, the Vice President of Colombia, at 9 July 2001 United Nations Conference on the Illicit Trade in Small Arms and Light Weapons in All Its Aspects, quoted the Millennium Report of the UN Secretary-General to the General Assembly, in which Kofi Annan said that small arms could be described as WMD because the fatalities they cause "dwarf that of all other weapons systems – and in most years greatly exceed the toll of the atomic bombs that devastated Hiroshima and Nagasaki".
An additional condition often implicitly applied to WMD is that the use of the weapons must be strategic. In other words, they would be designed to "have consequences far outweighing the size and effectiveness of the weapons themselves". The strategic nature of WMD also defines their function in the military doctrine of total war as targeting the means a country would use to support and supply its war effort, specifically its population, industry, and natural resources.
Within U.S. civil defense organizations, the category is now Chemical, Biological, Radiological, Nuclear, and Explosive (CBRNE), which defines WMD as:
(1) Any explosive, incendiary, poison gas, bomb, grenade, or rocket having a propellant charge of more than four ounces [113 g], missile having an explosive or incendiary charge of more than one-quarter ounce [7 g], or mine or device similar to the above. (2) Poison gas. (3) Any weapon involving a disease organism. (4) Any weapon that is designed to release radiation at a level dangerous to human life.
Military definition
For the general purposes of national defense, the U.S. Code defines a weapon of mass destruction as:
any weapon or device that is intended, or has the capability, to cause death or serious bodily injury to a significant number of people through the release, dissemination, or impact of:
toxic or poisonous chemicals or their precursors
a disease organism
radiation or radioactivity
For the purposes of the prevention of weapons proliferation, the U.S. Code defines weapons of mass destruction as "chemical, biological, and nuclear weapons, and chemical, biological, and nuclear materials used in the manufacture of such weapons".
Criminal (civilian) definition
For the purposes of U.S. criminal law concerning terrorism, weapons of mass destruction are defined as:
any "destructive device" defined as any explosive, incendiary, or poison gas – bomb, grenade, rocket having a propellant charge of more than four ounces, missile having an explosive or incendiary charge of more than one-quarter ounce, mine, or device similar to any of the devices described in the preceding clauses
any weapon that is designed or intended to cause death or serious bodily injury through the release, dissemination, or impact of toxic or poisonous chemicals, or their precursors
any weapon involving a biological agent, toxin, or vector
any weapon that is designed to release radiation or radioactivity at a level dangerous to human life
The Federal Bureau of Investigation's definition is similar to that presented above from the terrorism statute:
any "destructive device" as defined in Title 18 USC Section 921: any explosive, incendiary, or poison gas – bomb, grenade, rocket having a propellant charge of more than four ounces, missile having an explosive or incendiary charge of more than one-quarter ounce, mine, or device similar to any of the devices described in the preceding clauses
any weapon designed or intended to cause death or serious bodily injury through the release, dissemination, or impact of toxic or poisonous chemicals or their precursors
any weapon involving a disease organism
any weapon designed to release radiation or radioactivity at a level dangerous to human life
any device or weapon designed or intended to cause death or serious bodily injury by causing a malfunction of or destruction of an aircraft or other vehicle that carries humans or of an aircraft or other vehicle whose malfunction or destruction may cause said aircraft or other vehicle to cause death or serious bodily injury to humans who may be within range of the vector in its course of travel or the travel of its debris.
Indictments and convictions for possession and use of WMD such as truck bombs, pipe bombs, shoe bombs, and cactus needles coated with a biological toxin have been obtained under 18 USC 2332a.
As defined by 18 USC §2332 (a), a Weapon of Mass Destruction is:
(A) any destructive device as defined in section 921 of the title;
(B) any weapon that is designed or intended to cause death or serious bodily injury through the release, dissemination, or impact of toxic or poisonous chemicals, or their precursors;
(C) any weapon involving a biological agent, toxin, or vector (as those terms are defined in section 178 of this title); or
(D) any weapon that is designed to release radiation or radioactivity at a level dangerous to human life;
Under the same statute, conspiring, attempting, threatening, or using a Weapon of Mass Destruction may be imprisoned for any term of years or for life, and if resulting in death, be punishable by death or by imprisonment for any terms of years or for life. They can also be asked to pay a maximum fine of $250,000.
The Washington Post reported on 30 March 2006: "Jurors asked the judge in the death penalty trial of Zacarias Moussaoui today to define the term 'weapons of mass destruction' and were told it includes airplanes used as missiles". Moussaoui was indicted and tried for conspiracy to both destroy aircraft and use weapons of mass destruction, among others.
The surviving Boston Marathon bombing perpetrator, Dzhokhar Tsarnaev, was charged in June 2013 with the federal offense of "use of a weapon of mass destruction" after he and his brother Tamerlan Tsarnaev allegedly placed crude shrapnel bombs, made from pressure cookers packed with ball bearings and nails, near the finish line of the Boston Marathon. He was convicted in April 2015. The bombing resulted in three deaths and at least 264 injuries.
International law
The development and use of WMD is governed by several international conventions and treaties.
Use, possession, and access
Nuclear weapons
Nuclear weapons use the energy inside of an atom's nucleus to create massive explosions. This goal is achieved through nuclear fission and fusion.
Nuclear fission is when the nucleus of an atom is split into smaller nuclei. This process can be induced by shooting a neutron at the nucleus of an atom. When the neutron is absorbed by the atom, it becomes unstable, causing it to split and release energy. Modern nuclear weapons start this process by detonating chemical explosives around a pit of either uranium-235 or plutonium-239 metal. The force from this detonation is directed inwards, causing the pit of uranium or plutonium to compress to a dense point. Once the uranium/plutonium is dense enough, neutrons are then injected. This starts a fission chain reaction also known as an atomic explosion.
Nuclear fusion is essentially the opposite of fission. It is the fusing together of nuclei, not the splitting of it. When exposed to extreme pressure and temperature, some lightweight nuclei can fuse together and form heavier nuclei, releasing energy in the process. Fusion weapons (also known as “thermonuclear” or “hydrogen” weapons) use the fission process to initiate fusion. Fusion weapons use the energy released from a fission explosion to fuse hydrogen isotopes together. The energy released from these weapons creates a fireball, which reaches tens of million degrees. A temperature of this magnitude is similar to the temperature found at center of the sun, so it shouldn't be any surprise to learn that the sun runs on fusion as well.
The only country to have used a nuclear weapon in war is the United States, which dropped two atomic bombs on the Japanese cities of Hiroshima and Nagasaki during World War II.
At the start of 2023, nine states—the United States, Russia, the United Kingdom, France, China, India, Pakistan, the Democratic People's Republic of Korea (DPRK, or North Korea) and Israel—together possessed approximately 12 512 nuclear weapons, of which 9576 were considered to be potentially operationally available. An estimated 3844 of these warheads were deployed with operational forces, including about 2000 that were kept in a state of high operational alert—the same number as the previous year.
South Africa developed a small nuclear arsenal in the 1980s but disassembled them in the early 1990s, making it the only country to have fully given up an independently developed nuclear weapons arsenal. Belarus, Kazakhstan, and Ukraine inherited stockpiles of nuclear arms following the break-up of the Soviet Union, but relinquished them to the Russian Federation.
Countries where nuclear weapons are deployed through nuclear sharing agreements include Belgium, Germany, Italy, the Netherlands, and Turkey.
Biological weapons
The history of biological warfare goes back at least to the Mongol siege of Caffa in 1346 and possibly much farther back to antiquity. It is believed that the Ancient Greeks contaminated their adversaries' wells by placing animal corpses in them. However, only by the turn of the 20th century did advances in microbiology allow for the large-scale weaponization of pathogens. During First World War, German military attempted to introduce anthrax into Allied livestock. In Second World War, Japan conducted aerial attacks on China using fleas carrying the bubonic plague. During the 20th century, at least nine states have operated offensive biological weapons programs, including Canada (1946–1956), France (1921–1972), Iraq (1985–1990s), Japan (1930s–1945), Rhodesia, South Africa (1981–1993), the Soviet Union (1920s–1992), the United Kingdom (1934–1956), and the United States (1943–1969). The Japanese biological weapons program, which was run by the secret Imperial Japanese Army Unit 731 during the Sino-Japanese War (1937–1945), became infamous for conducting often fatal human experiments on prisoners and producing biological weapons for combat use. The Soviet Union covertly operated the world's largest, longest, and most sophisticated biological weapons program, in violation of its obligations under international law.
International restrictions on biological warfare began with the 1925 Geneva Protocol, which prohibits the use but not the possession or development of biological and chemical weapons. Upon ratification of the Geneva Protocol, several countries made reservations regarding its applicability and use in retaliation. Due to these reservations, it was in practice a "no-first-use" agreement only. The 1972 Biological Weapons Convention (BWC) supplements the Geneva Protocol by prohibiting the development, production, acquisition, transfer, stockpiling, and use of biological weapons. Having entered into force on 26 March 1975, the BWC was the first multilateral disarmament treaty to ban the production of an entire category of weapons of mass destruction. As of March 2021, 183 states have become party to the treaty.
Chemical weapons
Chemical weapons have been used around the world by various civilizations since ancient times. The oldest reported case of a chemical substance being used as a weapon was in 256 AD during the siege of Dura-Europos. A mixture of tar and sulfur was used to produce sulfur oxides, which helped take control of the city. In the industrial era, chemical weapons were used extensively by both sides during World War I, and by the Axis powers during World War II (both in battle and in extermination camp gas chambers) though Allied powers also stockpiled them.
International restrictions on chemical warfare began with the Hague Conventions of 1899 and 1907, and was expanded significantly by the 1925 Geneva Protocol. These treaties prohibited the use of poisons or chemical agents in international warfare, but did not place restrictions on development or weapon stockpiles. Since 1997, the Chemical Weapons Convention (CWC) has expanded restrictions to prohibit any use and development of chemical weapons except for very limited purposes (research, medical, pharmaceutical or protective). As of 2018, a handful of countries have known inventories, and many are in the process of being safely destroyed. Nonetheless, proliferation and use in war zones remains an active concern, most recently the use of chemical weapons in the Syrian Civil War.
Ethics and international legal status
Some commentators classify some or all the uses of nuclear, chemical, or biological weapons during wartime as a war crime (or crime against humanity if widespread) because they kill civilians (who are protected by the laws of war) indiscriminately or are specifically prohibited by international treaties (which have become more comprehensive over time). Proponents of use say that specific uses of such weapons have been necessary for defense or to avoid more deaths in a protracted war. The tactic of terror bombing from aircraft, and generally targeting cities with area bombardment or saturation carpet bombing has also been criticized, defended, and prohibited by treaty in the same way; the destructive effect of conventional saturation bombing is similar to that of a nuclear weapon.
United States politics
Due to the potentially indiscriminate effects of WMD, the fear of a WMD attack has shaped political policies and campaigns, fostered social movements, and has been the central theme of many films. Support for different levels of WMD development and control varies nationally and internationally. Yet understanding of the nature of the threats is not high, in part because of imprecise usage of the term by politicians and the media.
Fear of WMD, or of threats diminished by the possession of WMD, has long been used to catalyze public support for various WMD policies. They include mobilization of pro- and anti-WMD campaigners alike, and generation of popular political support. The term WMD may be used as a powerful buzzword or to generate a culture of fear. It is also used ambiguously, particularly by not distinguishing among the different types of WMD.
A television commercial called Daisy, promoting Democrat Lyndon Johnson's 1964 presidential candidacy, invoked the fear of a nuclear war and was an element in Johnson's subsequent election.
Later, United States' President George W. Bush used the threat of potential WMD in Iraq as justification for the 2003 invasion of Iraq. Broad reference to Iraqi WMD in general was seen as an element of President Bush's arguments. The claim that Iraq possessed Weapons of Mass Destruction (WMD) was a major factor that led to the invasion of Iraq in 2003 by Coalition forces.
Over 500 munitions containing mustard agent and sarin were discovered throughout Iraq since 2003; they were made in the 1980s and are no longer usable as originally intended due to corrosion.
The American Heritage Dictionary defines a weapon of mass destruction as: "a weapon that can cause widespread destruction or kill large numbers of people, especially a nuclear, chemical, or biological weapon." In other words, it does not have to be nuclear, biological or chemical (NBC). For example, Dzhokhar Tsarnaev, one of the perpetrators of the Boston Marathon bombing, was charged under United States law 18 U.S.C. 2332A for using a weapon of mass destruction and that was a pressure cooker bomb. In other words, it was a weapon that caused large-scale death and destruction, without being an NBC weapon.
Media coverage
In March 2004, the Center for International and Security Studies at Maryland (CISSM) released a report examining the media's coverage of WMD issues during three separate periods: nuclear weapons tests by India and Pakistan in May 1998; the U.S. announcement of evidence of a North Korean nuclear weapons program in October 2002; and revelations about Iran's nuclear program in May 2003. The CISSM report argues that poor coverage resulted less from political bias among the media than from tired journalistic conventions. The report's major findings were that:
In a separate study published in 2005, a group of researchers assessed the effects reports and retractions in the media had on people's memory regarding the search for WMD in Iraq during the 2003 Iraq War. The study focused on populations in two coalition countries (Australia and the United States) and one opposed to the war (Germany). Results showed that U.S. citizens generally did not correct initial misconceptions regarding WMD, even following disconfirmation; Australian and German citizens were more responsive to retractions. Dependence on the initial source of information led to a substantial minority of Americans exhibiting false memory that WMD were indeed discovered, while they were not. This led to three conclusions:
The repetition of tentative news stories, even if they are subsequently disconfirmed, can assist in the creation of false memories in a substantial proportion of people.
Once information is published, its subsequent correction does not alter people's beliefs unless they are suspicious about the motives underlying the events the news stories are about.
When people ignore corrections, they do so irrespective of how certain they are that the corrections occurred.
A poll conducted between June and September 2003 asked people whether they thought evidence of WMD had been discovered in Iraq since the war ended. They were also asked which media sources they relied upon. Those who obtained their news primarily from Fox News were three times as likely to believe that evidence of WMD had been discovered in Iraq than those who relied on PBS and NPR for their news, and one third more likely than those who primarily watched CBS.
Based on a series of polls taken from June–September 2003.
In 2006, Fox News reported the claims of two Republican lawmakers that WMDs had been found in Iraq, based upon unclassified portions of a report by the National Ground Intelligence Center. Quoting from the report, Senator Rick Santorum said "Since 2003, coalition forces have recovered approximately 500 weapons munitions which contain degraded mustard or sarin nerve agent". According to David Kay, who appeared before the U.S. House Armed Services Committee to discuss these badly corroded munitions, they were leftovers, many years old, improperly stored or destroyed by the Iraqis. Charles Duelfer agreed, stating on NPR's Talk of the Nation: "When I was running the ISG – the Iraq Survey Group – we had a couple of them that had been turned in to these IEDs, the improvised explosive devices. But they are local hazards. They are not a major, you know, weapon of mass destruction."
Later, wikileaks would show that WMDs of these kinds continued to be found as the Iraqi occupation continued.
Many news agencies, including Fox News, reported the conclusions of the CIA that, based upon the investigation of the Iraq Survey Group, WMDs are yet to be found in Iraq.
Public perceptions
Awareness and opinions of WMD have varied during the course of their history. Their threat is a source of unease, security, and pride to different people. The anti-WMD movement is embodied most in nuclear disarmament, and led to the formation of the British Campaign for Nuclear Disarmament in 1957.
In order to increase awareness of all kinds of WMD, in 2004 the nuclear physicist and Nobel Peace Prize winner Joseph Rotblat inspired the creation of The WMD Awareness Programme to provide trustworthy and up to date information on WMD worldwide.
In 1998, the University of New Mexico's Institute for Public Policy released their third report on U.S. perceptions – including the general public, politicians and scientists – of nuclear weapons since the breakup of the Soviet Union. Risks of nuclear conflict, proliferation, and terrorism were seen as substantial.
While maintenance of the U.S. nuclear arsenal was considered above average in importance, there was widespread support for a reduction in the stockpile, and very little support for developing and testing new nuclear weapons.
Also in 1998, nuclear weapons became an issue in India's election of March, in relation to political tensions with neighboring Pakistan. Prior to the election the Bharatiya Janata Party (BJP) announced it would "declare India a nuclear weapon state" after coming to power.
BJP won the elections, and on 14 May, three days after India tested nuclear weapons for the second time, a public opinion poll reported that a majority of Indians favored the country's nuclear build-up.
On 15 April 2004, the Program on International Policy Attitudes (PIPA) reported that U.S. citizens showed high levels of concern regarding WMD, and that preventing the spread of nuclear weapons should be "a very important U.S. foreign policy goal", accomplished through multilateral arms control rather than the use of military threats.
A majority also believed the United States should be more forthcoming with its biological research and its Nuclear Non-Proliferation Treaty commitment of nuclear arms reduction.
A Russian opinion poll conducted on 5 August 2005 indicated half the population believed new nuclear powers have the right to possess nuclear weapons. 39% believed the Russian stockpile should be reduced, though not eliminated.
In popular culture
Weapons of mass destruction and their related impacts have been a mainstay of popular culture since the beginning of the Cold War, as both political commentary and humorous outlet. The actual phrase "weapons of mass destruction" has been used similarly and as a way to characterise any powerful force or product since the Iraqi weapons crisis in the lead up to the Coalition invasion of Iraq in 2003. Science-fiction may introduce novel weapons of mass destruction with much greater yields or impact than anything in reality.
Common hazard symbols
Radioactive weaponry or hazard symbol
The international radioactivity symbol (also known as trefoil) first appeared in 1946, at the University of California, Berkeley Radiation Laboratory. At the time, it was rendered as magenta, and was set on a blue background.
It is drawn with a central circle of radius R, the blades having an internal radius of 1.5R and an external radius of 5R, and separated from each other by 60°. It is meant to represent a radiating atom.
The International Atomic Energy Agency found that the trefoil radiation symbol is unintuitive and can be variously interpreted by those uneducated in its meaning; therefore, its role as a hazard warning was compromised as it did not clearly indicate "danger" to many non-Westerners and children who encountered it. As a result of research, a new radiation hazard symbol (ISO 21482) was developed in 2007 to be placed near the most dangerous parts of radiation sources featuring a skull, someone running away, and using a red rather than yellow background.
The red background is intended to convey urgent danger, and the sign is intended to be used on equipment where very strong ionizing radiation can be encountered if the device is dismantled or otherwise tampered with. The intended use of the sign is not in a place where the normal user will see it, but in a place where it will be seen by someone who has started to dismantle a radiation-emitting device or equipment. The aim of the sign is to warn people such as scrap metal workers to stop work and leave the area.
Biological weaponry or hazard symbol
Developed by Dow Chemical company in the 1960s for their containment products.
According to Charles Dullin, an environmental-health engineer who contributed to its development:
| Technology | Weapons of mass destruction | null |
53179 | https://en.wikipedia.org/wiki/Boeing%20B-29%20Superfortress | Boeing B-29 Superfortress | The Boeing B-29 Superfortress was an American four-engined propeller-driven heavy bomber, designed by Boeing and flown primarily by the United States during World War II and the Korean War. Named in allusion to its predecessor, the Boeing B-17 Flying Fortress, the Superfortress was designed for high-altitude strategic bombing, but also excelled in low-altitude night incendiary bombing, and in dropping naval mines to blockade Japan. B-29s dropped the atomic bombs on Hiroshima and Nagasaki, the only aircraft ever to drop nuclear weapons in combat.
One of the largest aircraft of World War II, the B-29 was designed with state-of-the-art technology, which included a pressurized cabin, dual-wheeled tricycle landing gear, and an analog computer-controlled fire-control system that allowed one gunner and a fire-control officer to direct four remote machine gun turrets. The $3 billion cost of design and production (equivalent to $ billion in 2022), far exceeding the $1.9 billion cost of the Manhattan Project, made the B-29 program the most expensive of the war. The B-29 remained in service in various roles throughout the 1950s, being retired in the early 1960s after 3,970 had been built. A few were also used as flying television transmitters by the Stratovision company. The Royal Air Force flew the B-29 with the service name Washington from 1950 to 1954 when the jet-powered Canberra entered service.
The B-29 was the progenitor of a series of Boeing-built bombers, transports, tankers, reconnaissance aircraft, and trainers. For example, the re-engined B-50 Superfortress Lucky Lady II became the first aircraft to fly around the world non-stop, during a 94-hour flight in 1949. The Boeing C-97 Stratofreighter airlifter, which was first flown in 1944, was followed in 1947 by its commercial airliner variant, the Boeing Model 377 Stratocruiser. In 1948, Boeing introduced the KB-29 tanker, followed in 1950 by the Model 377-derivative KC-97. A line of outsized-cargo variants of the Stratocruiser is the GuppyMini GuppySuper Guppy, which remain in service with NASA and other operators. The Soviet Union produced 847 Tupolev Tu-4s, an unlicensed reverse-engineered copy of the B-29. Twenty-two B-29s have survived to preservation; while the majority are on static display at museums, 2 airframes, FIFI and Doc, still fly.
Design and development
Boeing began work on long-range bombers in 1938. Boeing's design study for the Model 334 was a pressurized derivative of the Boeing B-17 Flying Fortress with nosewheel undercarriage. Although the Air Corps lacked funds to pursue the design, Boeing continued development with its own funds as a private venture. In December 1939, the Air Corps issued a formal specification for a so-called "superbomber" that could deliver of bombs to a target away, and at a speed of . Boeing's previous private venture studies formed the starting point for its response to the Air Corps formal specification.
On 29 January 1940, the United States Army Air Corps issued a request to five major aircraft manufacturers to submit designs for a four-engine bomber with a range of . Boeing submitted its Model 345 on 11 May 1940, in competition with designs from Consolidated Aircraft (the Model 33, which later became the B-32), Lockheed (the Lockheed XB-30), and Douglas (the Douglas XB-31).
Douglas and Lockheed soon abandoned work on their projects, but Boeing received an order for two flying prototypes, which were given the designation XB-29, and an airframe for static testing on 24 August 1940, with the order being revised to add a third flying aircraft on 14 December. Consolidated continued to work on its Model 33, as it was seen by the Air Corps as a backup if there were problems with Boeing's design. These designs were evaluated, and on 6 September orders were placed for two experimental models each from Boeing and Consolidated Aircraft, which became the Boeing B-29 Superfortress and the Consolidated B-32 Dominator. These were known as very long range (VLR) bombers; the name "Superfortress" was not assigned until March 1944. On 17 May 1941, Boeing received an initial production order for 14 service test aircraft and 250 production bombers; this being increased to 500 aircraft in January 1942.
Manufacturing the B-29 was a complex task that involved four main-assembly factories. There were two Boeing operated plants at Renton, Washington (Boeing Renton Factory), and one in Wichita, Kansas (now Spirit AeroSystems), a Bell plant at Marietta, Georgia, near Atlanta ("Bell-Atlanta"), and a Martin plant at Bellevue, Nebraska ("Martin-Omaha" – Offutt Field). Thousands of subcontractors were also involved in the project. The first prototype made its maiden flight from Boeing Field, Seattle, on 21 September 1942. The combined effects of the aircraft's highly advanced design, challenging requirements, immense pressure for production, and hurried development caused setbacks. Unlike the unarmed first prototype, the second was fitted with a Sperry defensive armament system using remote-controlled gun turrets sighted by periscopes and first flew on 30 December 1942, although the flight was terminated due to a serious engine fire.
On 18 February 1943, the second prototype, flying out of Boeing Field in Seattle, experienced an engine fire and crashed. The crash killed Boeing test pilot Edmund T. Allen and his 10-man crew, 20 workers at the Frye Meat Packing Plant and a Seattle firefighter. Changes to the production craft came so often and so fast that, in early 1944, B-29s flew from the production lines directly to modification depots for extensive rebuilds to incorporate the latest changes. AAF-contracted modification centers and its own air depot system struggled to handle the scope of the requirements. Some facilities lacked hangars capable of housing the giant B-29, requiring outdoor work in freezing weather, further delaying necessary modification. By the end of 1943, although almost 100 aircraft had been delivered, only 15 were airworthy. This prompted an intervention by General Hap Arnold to resolve the problem, with production personnel being sent from the factories to the modification centers to speed availability of sufficient aircraft to equip the first bomb groups in what became known as the "Battle of Kansas". This resulted in 150 aircraft being modified in the five weeks, between 10 March and 15 April 1944.
The most common cause of maintenance headaches and catastrophic failures was the engines. Although the Wright R-3350 Duplex-Cyclone radial engines later became a trustworthy workhorse in large piston-engined aircraft, early models were beset with dangerous reliability problems. This problem was not fully cured until the aircraft was fitted with the more powerful Pratt & Whitney R-4360 "Wasp Major" in the B-29D/B-50 program, which arrived too late for World War II. Interim measures included cuffs placed on propeller blades to divert a greater flow of cooling air into the intakes, which had baffles installed to direct a stream of air onto the exhaust valves. Oil flow to the valves was also increased, asbestos baffles were installed around rubber push rod fittings to prevent oil loss, thorough pre-flight inspections were made to detect unseated valves, and mechanics frequently replaced the uppermost five cylinders (every 25 hours of engine time) and the entire engines (every 75 hours).
Pilots, including the present-day pilots of the Commemorative Air Force's Fifi, one of the last two remaining flying B-29s, describe flight after takeoff as being an urgent struggle for airspeed (generally, flight after takeoff should consist of striving for altitude). Radial engines need airflow to keep them cool, and failure to get up to speed as soon as possible could result in an engine failure and risk of fire. One useful technique was to check the magnetos while already on takeoff roll rather than during a conventional static engine-runup before takeoff.
The $3 billion cost of design and production (equivalent to $ billion in 2022), far exceeding the $1.9 billion cost of the Manhattan Project, made the B-29 program the most expensive of the war. Unit cost was US$639,188 (prototype cost $3,392,396.60)
Features
Defensive gun turret emplacements
In wartime, the B-29 was capable of flight at altitudes up to , at speeds of up to (true airspeed). This was its best defense because Japanese fighters could barely reach that altitude, and few could catch the B-29 even if they did attain that altitude.
The General Electric Central Fire Control system on the B-29 directed four remotely controlled turrets armed with two .50 Browning M2 machine guns each. All weapons were aimed optically, with targeting computed by analog electrical instrumentation. There were five interconnected sighting stations located in the nose and tail positions and three Plexiglas blisters in the central fuselage. Five General Electric analog computers (one dedicated to each sight) increased the weapons' accuracy by compensating for factors such as airspeed, lead, gravity, temperature and humidity. The computers also allowed a single gunner to operate two or more turrets (including tail guns) simultaneously. The gunner in the upper position acted as fire control officer, managing the distribution of turrets among the other gunners during combat. The tail position initially had two .50 Browning machine guns and a single M2 20 mm cannon. Later aircraft had the 20 mm cannon removed, sometimes replaced by a third machine gun.
In early 1945, Major General Curtis Lemay, commander of XXI Bomber Command—the Marianas-based B-29-equipped bombing force—ordered most of the defensive armament and remote-controlled sighting equipment removed from the B-29s under his command. The affected aircraft had the same reduced defensive firepower as the nuclear weapons-delivery intended Silverplate B-29 airframes and could carry greater fuel and bomb loads as a result of the change. The lighter defensive armament was made possible by a change in mission from high-altitude, daylight bombing with high explosive bombs to low-altitude night raids using incendiary bombs. As a consequence of that requirement, Bell Atlanta (BA) produced a series of 311 B-29Bs that had turrets and sighting equipment omitted, except for the tail position, which was fitted with AN/APG-15 fire-control radar. That version could also have an improved APQ-7 "Eagle" bombing-through-overcast radar fitted in an airfoil-shaped radome under the fuselage. Most of those aircraft were assigned to the 315th Bomb Wing, Northwest Field, Guam.
Pressurization
The crew would enjoy, for the first time in a bomber, full-pressurization comfort. This first-ever cabin pressure system for an Allied production bomber was developed for the B-29 by Garrett AiResearch. Both the forward and rear crew compartments were to be pressurized, but the designers had to decide whether to have bomb bays that were not pressurized or a fully pressurized fuselage that would have to be de-pressurized prior to opening the bomb bay doors. The solution was to have bomb bays that were not pressurized and a long tunnel joining the forward and rear crew compartments. Crews could use the tunnel if necessary to crawl from one pressurized compartment to the other.
Operational history
World War II
In September 1941, the United States Army Air Forces' plans for war against Germany and Japan proposed basing the B-29 in Egypt for operations against Germany, as British airbases were likely to be overcrowded. Air Force planning throughout 1942 and early 1943 continued to have the B-29 deployed initially against Germany, transferring to the Pacific only after the end of the war in Europe. By the end of 1943, plans had changed, partly due to production delays, and the B-29 was dedicated to the Pacific Theater. A new plan implemented at the direction of President Franklin D. Roosevelt as a promise to China, called Operation Matterhorn, deployed the B-29 units to attack Japan from four forward bases in southern China, with five main bases in India, and to attack other targets in the region from China and India as needed. The Chengdu region was eventually chosen over the Guilin region to avoid having to raise, equip, and train 50 Chinese divisions to protect the advanced bases from Japanese ground attack. The XX Bomber Command, initially intended to be two combat wings of four groups each, was reduced to a single wing of four groups because of the lack of availability of aircraft, automatically limiting the effectiveness of any attacks from China.
This was an extremely costly scheme, as there was no overland connection available between India and China, and all supplies had to be flown over the Himalayas, either by transport aircraft or by B-29s themselves, with some aircraft being stripped of armor and guns and used to deliver fuel. B-29s started to arrive in India in early April 1944. The first B-29 flight to airfields in China (over the Himalayas, or "The Hump") took place on 24 April 1944. The first B-29 combat mission was flown on 5 June 1944, with 77 out of 98 B-29s launched from India bombing the railroad shops in Bangkok and elsewhere in Thailand. Five B-29s were lost during the mission, none to hostile fire.
Forward base in China
On 5 June 1944, B-29s raided Bangkok, in what is reported as a test before being deployed against the Japanese home islands. Sources do not report from where they launched and vary as to the numbers involved—77, 98, and 114 being claimed. Targets were Bangkok's Memorial Bridge and a major power plant. Bombs fell over two kilometers away, damaged no civilian structures, but destroyed some tram lines, and destroyed both a Japanese military hospital and the Japanese secret police headquarters. On 15 June 1944, 68 B-29s took off from bases around Chengdu, 47 B-29s bombed the Imperial Iron and Steel Works at Yawata, Fukuoka Prefecture, Japan. This was the first attack on Japanese islands since the Doolittle raid in April 1942. The first B-29 combat losses occurred during this raid, with one B-29 destroyed on the ground by Japanese fighters after an emergency landing in China, one lost to anti-aircraft fire over Yawata, and another, the Stockett's Rocket (after Capt. Marvin M. Stockett, Aircraft Commander) B-29-1-BW 42-6261, disappeared after takeoff from Chakulia, India, over the Himalayas (12 KIA, 11 crew and one passenger). This raid, which did little damage to the target, with only one bomb striking the target factory complex, nearly exhausted fuel stocks at the Chengdu B-29 bases, resulting in a slow-down of operations until the fuel stockpiles could be replenished. Starting in July, the raids against Japan from Chinese airfields continued at relatively low intensity. Japan was bombed on:
7 July 1944 (14 B-29s)
29 July (70+)
10 August (24)
20 August (61)
8 September (90)
26 September (83)
25 October (59)
12 November (29)
21 November (61)
19 December (36)
6 January 1945 (49)
B-29s were withdrawn from airfields in China by the end of January 1945. Throughout the prior period, B-29 raids were also launched from China and India against many other targets throughout Southeast Asia, including a series of raids on Singapore and Thailand. On 2 November 1944, 55 B-29s raided Bangkok's Bang Sue marshaling yards in the largest raid of the war. Seven RTAF Nakajima Ki-43 Hayabusas from Foong Bin (Air Group) 16 and 14 IJAAF Ki-43s attempted intercept. RTAF Flt Lt Therdsak Worrasap attacked a B-29, damaging it, but was shot down by return fire. One B-29 was lost, possibly the one damaged by Flt Lt Therdsak. On 14 April 1945, a second B-29 raid on Bangkok destroyed two key power plants and was the last major attack conducted against Thai targets. The B-29 effort was gradually shifted to the new bases in the Mariana Islands in the Central Pacific, with the last B-29 combat mission from India flown on 29 March 1945.
New Mariana Islands air bases
In addition to the logistical problems associated with operations from China, the B-29 could reach only a limited part of Japan while flying from Chinese bases. The solution to this problem was to capture the Mariana Islands, which would bring targets such as Tokyo, about north of the Marianas within range of B-29 attacks. The Joint Chiefs of Staff agreed in December 1943 to seize the Marianas.
US forces invaded Saipan on 15 June 1944. Despite a Japanese naval counterattack which led to the Battle of the Philippine Sea and heavy fighting on land, Saipan was secured by 9 July. Operations followed against Guam and Tinian, with all three islands secured by August.
Naval construction battalions (Seabees) began at once to construct air bases suitable for the B-29, commencing even before the end of ground fighting. In all, five major airfields were built: two on the flat island of Tinian, one on Saipan, and two on Guam. Each was large enough to eventually accommodate a bomb wing consisting of four bomb groups, giving a total of 180 B-29s per airfield. These bases could be supplied by ship and, unlike the bases in China, were not vulnerable to attack by Japanese ground forces.
The bases became the launch sites for the large B-29 raids against Japan in the final year of the war. The first B-29 arrived on Saipan on 12 October 1944, and the first combat mission was launched from there on 28 October 1944, with 14 B-29s attacking the Truk atoll. The 73rd Bomb Wing launched the first mission against Japan from bases in the Marianas, on 24 November 1944, sending 111 B-29s to attack Tokyo. For this first attack on the Japanese capital since the Doolittle Raid in April 1942, 73rd Bomb Wing wing commander Brigadier General Emmett O'Donnell Jr. acted as mission command pilot in B-29 Dauntless Dotty.
The campaign of incendiary raids started with the bombardment of Kobe on 4 February 1945, then peaked early with the most destructive bombing raid in history (even when the later Silverplate-flown nuclear attacks on Hiroshima and Nagasaki are considered) on the night of 9–10 March 1945 on Tokyo. From then on, the raids intensified, being launched regularly until the end of the war. The attacks succeeded in devastating most large Japanese cities (with the exception of Kyoto and four that were reserved for nuclear attacks), and gravely damaged Japan's war industries. Although less publicly appreciated, the mining of Japanese ports and shipping routes (Operation Starvation) carried out by B-29s from April 1945 reduced Japan's ability to support its population and move its troops.
Nuclear weapons
The most famous B-29s were the Silverplate series, being extensively modified to carry nuclear weapons. Early consideration was given to using the British Lancaster as a nuclear bomber, as this would require less modification. However, the superior range and high-altitude performance of the B-29 made it a much better choice, and after the B-29 began to be modified in November 1943 for carrying the atomic bomb, the suggestion for using the Lancaster never came up again.
The most significant modification was the enlargement of the bomb bay enabling each aircraft to carry either the Thinman or Fatman weapons. These Silverplate bombers differed from other B-29s then in service by having fuel injection and reversible props. Also, to make a lighter aircraft, the Silverplate B-29s were stripped of all guns, except for those on the tail. Pilot Charles Sweeney credits the reversible props for saving Bockscar after making an emergency landing on Okinawa following the Nagasaki bombing.
Enola Gay, flown by Colonel Paul Tibbets, dropped the first bomb, called Little Boy, on Hiroshima on 6 August 1945. Enola Gay is fully restored and on display at the Smithsonian's Steven F. Udvar-Hazy Center, outside Dulles Airport near Washington, D.C. Bockscar, piloted by Major Charles W. Sweeney, dropped the second bomb, called Fat Man, on Nagasaki three days later. Bockscar is on display at the National Museum of the United States Air Force.
Following the surrender of Japan, called V-J Day, B-29s were used for other purposes. A number supplied POWs with food and other necessities by dropping barrels of rations on Japanese POW camps. In September 1945, a long-distance flight was undertaken for public relations purposes: Generals Barney M. Giles, Curtis LeMay, and Emmett O'Donnell Jr. piloted three specially modified B-29s from Chitose Air Base in Hokkaidō to Chicago Municipal Airport, continuing to Washington, D.C., the farthest nonstop distance () to that date flown by U.S. Army Air Forces aircraft and the first-ever nonstop flight from Japan to Chicago. Two months later, Colonel Clarence S. Irvine commanded another modified B-29, Pacusan Dreamboat, in a world-record-breaking long-distance flight from Guam to Washington, D.C., traveling in 35 hours, with a gross takeoff weight of . Almost a year later, in October 1946, the same B-29 flew nonstop from Oahu, Hawaii, to Cairo, Egypt, in less than 40 hours, demonstrating the possibility of routing airlines over the polar ice cap.
B-29s in Europe and Australia
Although considered for other theaters, and briefly evaluated in the UK, the B-29 was exclusively used in World War II in the Pacific Theatre. The use of YB-29-BW 41-36393, the so-named Hobo Queen, one of the service test aircraft flown around several British airfields in early 1944, was part of a "disinformation" program from its mention in an American-published Sternenbanner German-language propaganda leaflet from Leap Year Day in 1944, meant to be circulated within the Reich, with the intent to deceive the Germans into believing that the B-29 would be deployed to Europe.
American post-war military assistance programs loaned the RAF 87 Superfortresses, to equip eight RAF Bomber Command squadrons. The aircraft was known as the Washington B.1 in RAF service and served from March 1950 until the last bombers were returned in March 1954. Deployment was restricted to long-range training for strategic attacks against the Soviet Union, which was beyond the range of the RAF's Avro Lincolns. The phase-out was occasioned by deliveries of the English Electric Canberra bombers.
Three Washingtons modified for ELINT duties and a standard bomber version used for support by No. 192 Squadron RAF were decommissioned in 1958, being replaced by de Havilland Comet aircraft.
Two British Washington B.1 aircraft were transferred to the Royal Australian Air Force (RAAF) in 1952. They were attached to the Aircraft Research and Development Unit and used in trials conducted on behalf of the British Ministry of Supply. Both aircraft were placed in storage in 1956 and were sold for scrap in 1957.
Soviet Tupolev Tu-4
At the end of WWII, Soviet development of modern four-engine heavy bombers lagged behind the West. The Petlyakov Pe-8—the sole heavy bomber operated by the Soviet Air Forces—first flew in 1936. Intended to replace the obsolete Tupolev TB-3, only 93 Pe-8s were built by the end of WWII. During 1944 and 1945, four B-29s made emergency landings in Soviet territory after bombing raids on Japanese Manchuria and Japan. In accordance with Soviet neutrality in the Pacific War, the bombers were interned by the Soviets despite American requests for their return. Rather than return the aircraft, the Soviets reverse engineered the American B-29s and used them as a pattern for the Tupolev Tu-4.
On 31 July 1944, Ramp Tramp (serial number 42-6256), of the United States Army Air Forces 462nd (Very Heavy) Bomb Group was diverted to Vladivostok, Russia, after an engine failed and the propeller could not be feathered. This B-29 was part of a 100-aircraft raid against the Japanese Showa steel mill in Anshan, Manchuria. On 20 August 1944, Cait Paomat (42-93829), flying from Chengdu, was damaged by anti-aircraft gunfire during a raid on the Yawata Iron Works. Due to the damage it sustained, the crew elected to divert to the Soviet Union. The aircraft crashed in the foothills of Sikhote-Alin mountain range east of Khabarovsk after the crew bailed out.
On 11 November 1944, during a night raid on Omura in Kyushu, Japan, the General H. H. Arnold Special (42-6365) was damaged and forced to divert to Vladivostok in the Soviet Union. The crew was interned. On 21 November 1944, Ding How (42-6358) was damaged during a raid on an aircraft factory at Omura and was also forced to divert to Vladivostok.
The interned crews of these four B-29s were allowed to escape into American-occupied Iran in January 1945, but none of the B-29s were returned after Stalin ordered the Tupolev OKB to examine and copy the B-29 and produce a design ready for quantity production as soon as possible.
Because aluminum in the USSR was supplied in different gauges from that available in the US (metric vs imperial), the entire aircraft had to be extensively re-engineered. In addition, Tupolev substituted his own favored airfoil sections for those used by Boeing, with the Soviets themselves already having their own Wright R-1820-derived 18 cylinder radial engine, the Shvetsov ASh-73 of comparable power and displacement to the B-29's Duplex Cyclone radials available to power their design. In 1947, the Soviets debuted both the Tupolev Tu-4 (NATO ASCC code named Bull), and the Tupolev Tu-70 transport variant. The Soviets used tail-gunner positions similar to the B-29 in many later bombers and transports.
Transition to USAF
Production of the B-29 was phased out after WWII, with the last example completed by Boeing's Renton factory on 28 May 1946. Many aircraft went into storage, being declared excess inventory, and were ultimately scrapped as surplus. Others remained in the active inventory and equipped the Strategic Air Command when it formed on 21 March 1946. In particular, the "Silverplate" modified aircraft of the 509th Composite Group remained the only aircraft capable of delivering the atomic bomb, and so the unit was involved in the Operation Crossroads series of tests, with B-29 Dave's Dream dropping a Fat Man bomb in Test Able on 1 July 1946.
Some B-29s, fitted with filtered air sampling scoops, were used to monitor above-ground nuclear weapons testing by the US and the USSR by sampling airborne radioactive contamination. The USAF also used the aircraft for long-range weather reconnaissance (WB-29), for signals intelligence gathering (EB-29) and photographic reconnaissance (RB-29).
Korean War and postwar service
The B-29 was used in 1950–1953 in the Korean War. At first, the bomber was used in normal strategic day-bombing missions, although North Korea's few strategic targets and industries were quickly destroyed. More importantly, in 1950 numbers of Soviet MiG-15 jet fighters appeared over Korea, and after the loss of 28 aircraft, future B-29 raids were restricted to night missions, largely in a supply-interdiction role.
The B-29 dropped the VB-3 "Razon" (a range-controllable version of the earlier Azon guided ordnance device) and the VB-13 "Tarzon" MCLOS radio-controlled bombs in Korea, mostly for demolishing major bridges, like the ones across the Yalu River, and for attacks on dams. The aircraft also was used for numerous leaflet drops in North Korea, such as those for Operation Moolah.
A Superfortress of the 91st Strategic Reconnaissance Squadron flew the last B-29 mission of the war on 27 July 1953.
Over the course of the war, B-29s flew 20,000 sorties and dropped 200,000 tonnes (220,000 tons) of bombs. B-29 gunners were credited with shooting down 27 enemy aircraft. In turn 78 B-29s were lost; 57 B-29 and reconnaissance variants were lost in action and 21 were non-combat losses.
Soviet records show that one MiG-15 jet fighter was shot down by a B-29 during the war. This occurred on 6 December 1950, when a B-29 shot down Lieutenant N. Serikov.
With the arrival of the mammoth Convair B-36, the B-29 was reclassified as a medium bomber by the Air Force. The later B-50 Superfortress variant (initially designated B-29D) was able to handle auxiliary roles such as air-sea rescue, electronic intelligence gathering, air-to-air refueling, and weather reconnaissance.
The B-50D was replaced in its primary role during the early 1950s by the Boeing B-47 Stratojet, which in turn was replaced by the Boeing B-52 Stratofortress. The final active-duty KB-50 and WB-50 variants were phased out in the mid-1960s, with the final example retired in 1965. A total of 3,970 B-29s were built.
Variants
The variants of the B-29 were outwardly similar in appearance but were built around different wing center sections that affected the wingspan dimensions. The wing of the Renton-built B-29A-BN used a different subassembly process and was a foot longer in span. The Georgia-built B-29B-BA weighed less through armament reduction. A planned C series with more reliable R-3350s was not built.
Moreover, engine packages changed, including the type of propellers and range of the variable pitch. A notable example was the eventual 65 airframes (up to 1947's end) for the Silverplate and successor-name "Saddletree" specifications built for the Manhattan Project with Curtiss Electric reversible pitch propellers.
The other differences came through added equipment for varied mission roles. These roles included cargo carriers (CB); rescue aircraft (SB); weather ships (WB); and trainers (TB); and aerial tankers (KB).
Some were used for odd purposes such as flying relay television transmitters under the name of Stratovision.
The B-29D led progressively to the XB-44, and the family of B-50 Superfortress (which was powered by four Pratt & Whitney R-4360-35 Wasp Major engines).
Another role was as a mothership. This included being rigged for carrying the experimental parasite fighter aircraft, such as the McDonnell XF-85 Goblin and Republic F-84 Thunderjets as in flight lock on and offs. It was also used to develop the Airborne Early Warning program; it was the ancestor of various modern radar picket aircraft. A B-29 with the original Wright Duplex Cyclone powerplants was used to air-launch the Bell X-1 supersonic research rocket aircraft, as well as Cherokee rockets for the testing of ejection seats.
Some B-29s were modified to act as testbeds for various new systems or special conditions, including fire-control systems, cold-weather operations, and various armament configurations. Several converted B-29s were used to experiment with aerial refueling and re-designated as KB-29s. Perhaps the most important tests were conducted by the XB-29G. It carried prototype jet engines in its bomb bay, and lowered them into the air stream to conduct measurements.
Operators
Royal Australian Air Force (two former RAF aircraft for trials)
Royal Air Force (87 loaned from the USAF as the Washington B.1)
United States Army Air Forces
United States Air Force
United States Navy (four former USAF aircraft designated as P2B patrol bombers)
Soviet Air Forces (three USAAF B-29s made emergency landings in the USSR during WWII, and were never returned; they were reverse-engineered to make the Soviet Tupolev Tu-4 "Bull" bomber.)
Surviving aircraft
Twenty-two B-29s are preserved at various museums worldwide, including two flying examples; FIFI, which belongs to the Commemorative Air Force, and Doc, which belongs to Doc's Friends. Doc made its first flight in 60 years from Wichita, Kansas, on 17 July 2016. The public is being invited to inspect and take a short paid flight in Doc and Fifi at various venues.
Three of the Silverplate B-29s modified to drop nuclear bombs survived. Superfortress 44-86292 Enola Gay (nose number 82), which dropped the first atomic bomb, was fully restored and placed on display at the Smithsonian's Steven F. Udvar-Hazy Center of the National Air & Space Museum near Washington Dulles International Airport in 2003. The B-29 that dropped Fat Man on Nagasaki, Superfortress 44-27297 Bockscar (nose number 77), is restored and on display at the National Museum of the United States Air Force at Wright-Patterson AFB in Dayton, Ohio, posed with a replica of the Mark 3 Fat Man nuclear bomb. The third is Superfortress 45-21748, which was delivered on 9 August 1945 and is on display at the National Museum of Nuclear Science and History in Albuquerque, New Mexico.
Only two of the twenty-two museum aircraft are outside the United States: It's Hawg Wild at the Imperial War Museum Duxford and another at the KAI Aerospace Museum in Sachon, South Korea.
Accidents and incidents
Notable accidents and incidents involving B-29s include:
The 1947 crash of the Kee Bird in Greenland during a flight to the geographic North Pole, and its subsequent destruction in 1995 during a recovery attempt.
The 1948 Waycross B-29 crash, which resulted in the United States v. Reynolds lawsuit regarding state secrets privilege.
The 1948 Lake Mead Boeing B-29 crash during the "Sun Tracker" project that aimed to develop an intercontinental ballistic missile guidance system that used the sun for direction and positioning.
The 3 November 1948 crash at Bleaklow moor near Glossop, Derbyshire, England. Much of the wreckage is still exposed and can be reached by a walk from the summit of Snake Pass, starting along the Pennine Way footpath through Devil's Dyke.
On 11 April 1950 a B-29 departed Kirtland Air Force Base and crashed into a mountain on Manzano Base approximately three minutes later, killing the crew. Detonators were installed in the nuclear bomb on the aircraft. The bomb case was demolished and some high-explosive (HE) material burned in the fire. Both the weapon and the capsule of nuclear material were on board but the capsule was not inserted in the bomb for safety reasons, so no nuclear detonation was possible.
On 5 August 1950, a B-29 carrying a Mark 4 nuclear bomb crashed shortly after takeoff from Fairfield-Suisun Air Force Base with 20 men on board. Twelve men were killed in the crash, including Brigadier General Robert F. Travis, and another seven on the ground when the aircraft exploded. The base was later renamed after Travis.
Specifications
Notable appearances in media
| Technology | Specific aircraft | null |
53217 | https://en.wikipedia.org/wiki/Exotic%20atom | Exotic atom | An exotic atom is an otherwise normal atom in which one or more sub-atomic particles have been replaced by other particles. For example, electrons may be replaced by other negatively charged particles such as muons (muonic atoms) or pions (pionic atoms). Because these substitute particles are usually unstable, exotic atoms typically have very short lifetimes and no exotic atom observed so far can persist under normal conditions.
Muonic atoms
In a muonic atom (previously called a mu-mesic atom, now known to be a misnomer as muons are not mesons), an electron is replaced by a muon, which, like the electron, is a lepton. Since leptons are only sensitive to weak, electromagnetic and gravitational forces, muonic atoms are governed to very high precision by the electromagnetic interaction.
Since a muon is more massive than an electron, the Bohr orbits are closer to the nucleus in a muonic atom than in an ordinary atom, and corrections due to quantum electrodynamics are more important. Study of muonic atoms' energy levels as well as transition rates from excited states to the ground state therefore provide experimental tests of quantum electrodynamics.
Muon-catalyzed fusion is a technical application of muonic atoms.
Other muonic atoms can be formed when negative muons interact with ordinary matter. The muon in muonic atoms can either decay or get captured by a proton. Muon capture is very important in heavier muonic atoms, but shortens the muon's lifetime from 2.2 μs to only 0.08 μs.
Muonic hydrogen
Muonic hydrogen is like normal hydrogen with the electron replaced by a negative muon—that is a proton orbited by a muon. It is important in addressing the proton radius puzzle.
Muonic helium (Hydrogen-4.1)
The symbol 4.1H (Hydrogen-4.1) has been used to describe the exotic atom muonic helium (4He-μ), which is like helium-4 in having two protons and two neutrons. However one of its electrons is replaced by a muon, which also has charge –1. Because the muon's orbital radius is less than the electron's orbital radius (due to the mass ratio), the muon can be considered as a part of the nucleus. The atom then has a nucleus with two protons, two neutrons and one muon, with total nuclear charge +1 (from two protons and one muon) and only one electron outside, so that it is effectively an isotope of hydrogen instead of an isotope of helium. A muon's weight is approximately 0.1 Da so the isotopic mass is 4.1. Since there is only one electron outside the nucleus, the hydrogen-4.1 atom can react with other atoms. Its chemical behavior behaves more like a hydrogen atom than an inert helium atom.
Hadronic atoms
A hadronic atom is an atom in which one or more of the orbital electrons are replaced by a negatively charged hadron. Possible hadrons include mesons such as the pion or kaon, yielding a pionic atom or a kaonic atom (see Kaonic hydrogen), collectively called mesonic atoms; antiprotons, yielding an antiprotonic atom; and the particle, yielding a or sigmaonic atom.
Unlike leptons, hadrons can interact via the strong force, so the orbitals of hadronic atoms are influenced by nuclear forces between the nucleus and the hadron. Since the strong force is a short-range interaction, these effects are strongest if the atomic orbital involved is close to the nucleus, when the energy levels involved may broaden or disappear because of the absorption of the hadron by the nucleus. Hadronic atoms, such as pionic hydrogen and kaonic hydrogen, thus provide experimental probes of the theory of strong interactions, quantum chromodynamics.
Onium
An onium (plural: onia) is the bound state of a particle and its antiparticle. The classic onium is positronium, which consists of an electron and a positron bound together as a metastable state, with a relatively long lifetime of 142 ns in the triplet state. Positronium has been studied since the 1950s to understand bound states in quantum field theory. A recent development called non-relativistic quantum electrodynamics (NRQED) used this system as a proving ground.
Pionium, a bound state of two oppositely charged pions, is useful for exploring the strong interaction. This should also be true of protonium, which is a proton–antiproton bound state. Understanding bound states of pionium and protonium is important in order to clarify notions related to exotic hadrons such as mesonic molecules and pentaquark states. Kaonium, which is a bound state of two oppositely charged kaons, has not been observed experimentally yet.
The true analogs of positronium in the theory of strong interactions, however, are not exotic atoms but certain mesons, the quarkonium states, which are made of a heavy quark such as the charm or bottom quark and its antiquark. (Top quarks are so heavy that they decay through the weak force before they can form bound states.) Exploration of these states through non-relativistic quantum chromodynamics (NRQCD) and lattice QCD are increasingly important tests of quantum chromodynamics.
Muonium, despite its name, is not an onium state containing a muon and an antimuon, because IUPAC assigned that name to the system of an antimuon bound with an electron. However, the production of a muon–antimuon bound state, which is an onium (called true muonium), has been theorized. The same applies to the ditauonium (or "true tauonium") exotic QED atom.
Hypernuclear atoms
Atoms may be composed of electrons orbiting a hypernucleus that includes strange particles called hyperons. Such hypernuclear atoms are generally studied for their nuclear behaviour, falling into the realm of nuclear physics rather than atomic physics.
Quasiparticle atoms
In condensed matter systems, specifically in some semiconductors, there are states called excitons, which are bound states of an electron and an electron hole.
Exotic molecules
An exotic molecule contains one or more exotic atoms.
Di-positronium, two bound positronium atoms
Positronium hydride, a positronium atom bound to a hydrogen atom
"Exotic molecule" can also refer to a molecule having some other uncommon property such as pyramidal hexamethylbenzene and a Rydberg atom.
| Physical sciences | Atomic physics | Physics |
53227 | https://en.wikipedia.org/wiki/Vein | Vein | Veins () are blood vessels in the circulatory system of humans and most other animals that carry blood towards the heart. Most veins carry deoxygenated blood from the tissues back to the heart; exceptions are those of the pulmonary and fetal circulations which carry oxygenated blood to the heart. In the systemic circulation, arteries carry oxygenated blood away from the heart, and veins return deoxygenated blood to the heart, in the deep veins.
There are three sizes of veins: large, medium, and small. Smaller veins are called venules, and the smallest the post-capillary venules are microscopic that make up the veins of the microcirculation. Veins are often closer to the skin than arteries.
Veins have less smooth muscle and connective tissue and wider internal diameters than arteries. Because of their thinner walls and wider lumens they are able to expand and hold more blood. This greater capacity gives them the term of capacitance vessels. At any time, nearly 70% of the total volume of blood in the human body is in the veins. In medium and large sized veins the flow of blood is maintained by one-way (unidirectional) venous valves to prevent backflow. In the lower limbs this is also aided by muscle pumps, also known as venous pumps that exert pressure on intramuscular veins when they contract and drive blood back to the heart.
Structure
There are three sizes of vein, large, medium, and small. Smaller veins are called venules. The smallest veins are the post-capillary venules. Veins have a similar three-layered structure to arteries. The layers known as tunicae have a concentric arrangement that forms the wall of the vessel. The outer layer, is a thick layer of connective tissue called the tunica externa or adventitia; this layer is absent in the post-capillary venules. The middle layer, consists of bands of smooth muscle and is known as the tunica media. The inner layer, is a thin lining of endothelium known as the tunica intima. The tunica media in the veins is much thinner than that in the arteries as the veins are not subject to the high systolic pressures that the arteries are. There are valves present in many veins that maintain unidirectional flow.
Unlike arteries, the precise location of veins varies among individuals.
Veins close to the surface of the skin appear blue for a variety of reasons. The factors that contribute to this alteration of color perception are related to the light-scattering properties of the skin and the processing of visual input by the visual cortex, rather than the actual colour of the venous blood which is dark red.
Venous system
The venous system is the system of veins in the systemic and pulmonary circulations that return blood to the heart. In the systemic circulation the return is of deoxygenated blood from the organs and tissues of the body, and in the pulmonary circulation the pulmonary veins return oxygenated blood from the lungs to the heart. Almost 70% of the blood in the body is in the veins, and almost 75% of this blood is in the small veins and venules. All of the systemic veins are tributaries of the largest veins, the superior and inferior vena cava, which empty the oxygen-depleted blood into the right atrium of the heart. The thin walls of the veins, and their greater internal diameters (lumens) enable them to hold a greater volume of blood, and this greater capacitance gives them the term of capacitance vessels. This characteristic also allows for the accommodation of pressure changes in the system. The whole of the venous system, bar the post-capillary venules is a large volume, low pressure system. The venous system is often asymmetric, and whilst the main veins hold a relatively constant position, unlike arteries, the precise location of veins varies among individuals.
Veins vary in size from the smallest post-capillary venules, and more muscular venules, to small veins, medium veins, and large veins. The thickness of the walls of the veins varies as to their location – in the legs the vein walls are much thicker than those in the arms. In the circulatory system, blood first enters the venous system from capillary beds where arterial blood changes to venous blood.
Large arteries such as the thoracic aorta, subclavian, femoral and popliteal arteries lie close to a single vein that drains the same region. Other arteries are often accompanied by a pair of veins held in a connective tissue sheath. The accompanying veins are known as venae comitantes, or satellite veins, and they run on either side of the artery. When an associated nerve is also enclosed, the sheath is known as a neurovascular bundle. This close proximity of the artery to the veins helps in venous return due to the pulsations in the artery. It also allows for the promotion of heat transfer from the larger arteries to the veins in a counterflow exchange that helps to preserve normal body heat.
Venules
The first entry of venous blood is from the convergence of two or more capillaries into a microscopic, post-capillary venule. Post-capillary venules have a diameter of between 10 and 30 micrometres (μm), and are part of the microcirculation. Their endothelium is made up of flattened oval or polygon shaped cells surrounded by a basal lamina. Post-capillary venules are too small to have a smooth muscle layer and are instead supported by pericytes that wrap around them. Post-capillary venules become muscular venules when they reach a diameter of 50 μm, and can reach a diameter of 1 mm. These larger venules feed into small veins.
Small, medium, and large veins
The small veins merge to feed as tributaries into medium-sized veins. The medium veins feed into the large veins which include the internal jugular, and renal veins, and the venae cavae that carry the blood directly into the heart. The venae cavae enter the right atrium of the heart from above and below. From above, the superior vena cava carries blood from the arms, head, and chest to the right atrium of the heart, and from below, the inferior vena cava carries blood from the legs and abdomen to the right atrium. The inferior vena cava is the larger of the two. The inferior vena cava is retroperitoneal and runs to the right and roughly parallel to the abdominal aorta along the spine.
Deep, superficial, and perforator veins
The three main compartments of the venous system are the deep veins, the superficial veins, and the perforator veins. Superficial veins are those closer to the surface of the body, and have no corresponding arteries. Deep veins are deeper in the body and have corresponding arteries. Perforator veins drain from the superficial to the deep veins. These are usually referred to in the lower limbs and feet. Superficial veins include the very small spider veins of between 0.5 and 1 mm diameter, and reticular or feeder veins.
Venous plexuses
There are a number of venous plexuses where veins are grouped or sometimes combined in networks at certain body sites. The Batson venous plexus, runs through the inner vertebral column connecting the thoracic and pelvic veins. These veins are noted for being valveless, believed to be the reason for metastasis of certain cancers.
A subcutaneous venous plexus is continuous, and a high rate of flow is supplied by small arteriovenous anastomoses. The high rate of flow ensures heat transfer to the vein wall.
Venous valves
Blood flows back to the heart in the systemic deep veins, with the flow of blood maintained by one-way valves in the deep veins, superficial veins, and in the perforator veins. The venous valves serve to prevent regurgitation (backflow) due to the low pressure of veins, and the pull of gravity. They also serve to prevent the over-widening of the vein.
A venous valve is bicuspid (having two leaflets) and is formed by an infolding of part of the tunica intima on either side of the lumen of the veins. The leaflets are strengthened with collagen, and elastic fibres, and covered with endothelium. The endothelial cells on the surfaces of the leaflets facing the vein wall, are arranged transversely. On the leaflet surfaces that open to let the blood flow, the cells are arranged longitudinally in the direction of the flow. The leaflets are attached to the venous wall at their convex edges. Their margins are concave and are directed with the flow lying against the wall. As the valve forms, the vein wall where the leaflets attach, becomes dilated on each side. These widenings form the pockets, hollow cup-shaped regions, on the cardial side, known as the valvular sinuses. The endothelial cells in the sinuses are able to stretch twice as much as those in areas without valves. When the blood tries to reverse its direction (due to low venous pressure and the pull of gravity), the sinuses fill first closing the leaflets and keeping them together. Approximately 95% of the venous valves are in the small veins of less than 300 micrometres.
The deep veins of the lower limb include the common femoral vein, femoral vein, and the deep femoral vein; the popliteal vein, the tibial, and fibular veins. In the common femoral vein one valve is located above the saphenofemoral junction called the suprasaphenic valve. There are sometimes two valves in the same tract. In the femoral vein there are often three valves, the most constantly found valve is just below the joining of the deep femoral vein. The deep femoral vein and its perforators have valves. In the popliteal veins there are between one and three valves; in each posterior tibial vein there are between 8 and 19 valves, and in the anterior tibial veins there are between 8 and 11 valves.
In the superficial veins there are between one and seven valves along the thigh portion of the great saphenous vein (GSV); two to six below the knee and one to four in the marginal veins of the foot. There is a valve at the termination of the GSV known as the terminal valve to prevent reflux from the femoral vein A preterminal valve is located just below the openings of the tributaries to prevent reflux form these into the GSV. Incompetence of the GSV is a common cause of varicose veins.
The valves also divide the column of blood into segments which helps move the blood unidirectionally to the heart. Their action is supported by the action of skeletal muscle pumps that contract and compress the veins. A skeletal muscle is confined in its fascia and contraction of the muscle which makes it wider results In compression on the vein that pushes the blood forward. Valves in the perforating veins close when a calf muscle contracts, to prevent backflow from the deep veins to the superficial. There are more valves in the lower leg, due to increased gravitational pull, with the number decreasing as the veins travel to the hip. There are no valves in the veins of the thorax or abdomen.
There is a valve at the junction of the inferior vena cava (one of the great vessels) and the right atrium known as the valve of inferior vena cava also known as the eustachian valve. This valve is an embryological remnant and is insignificant in the adult. However, when persistent it can cause problems.
Circulatory routes
There are some separate parallel systemic circulatory routes that supply specific regions, and organs. They include the coronary circulation, the cerebral circulation, the bronchial circulation, and the renal circulation.
Coronary circulation
In the coronary circulation, the blood supply to the heart, is drained by cardiac veins (or coronary veins) that remove the deoxygenated blood from the heart muscle. These include the great cardiac vein, the middle cardiac vein, the small cardiac vein, the smallest cardiac veins, and the anterior cardiac veins. Cardiac veins carry blood with a poor level of oxygen, from the heart muscle to the right atrium. Most of the blood of the cardiac veins returns through the coronary sinus. The anatomy of the veins of the heart is very variable, but generally it is formed by the following veins: heart veins that go into the coronary sinus: the great cardiac vein, the middle cardiac vein, the small cardiac vein, the posterior vein of the left ventricle, and the oblique vein of the left atrium (oblique vein of Marshall). Heart veins that go directly to the right atrium: the anterior cardiac veins, and the smallest cardiac veins (Thebesian veins).
Bronchial circulation
In the bronchial circulation that supplies blood to the lung tissues, bronchial veins drain venous blood from the large main bronchi into the azygous vein, and ultimately the right atrium. Venous blood from the bronchi inside the lungs drains into the pulmonary veins and empties into the left atrium; since this blood never went through a capillary bed it was never oxygenated and so provides a small amount of shunted deoxygenated blood into the systemic circulation.
Cerebral circulation
In the cerebral circulation supplying the cerebrum the venous drainage can be separated into two subdivisions: superficial and deep.
The superficial system is composed of dural venous sinuses, which have walls composed of dura mater as opposed to a traditional vein. The dural sinuses are therefore located on the surface of the cerebrum. The most prominent of these sinuses is the superior sagittal sinus which flows in the sagittal plane under the midline of the cerebral vault, posteriorly and inferiorly to the confluence of sinuses, where the superficial drainage joins with the sinus that primarily drains the deep venous system. From here, two transverse sinuses bifurcate and travel laterally and inferiorly in an S-shaped curve that forms the sigmoid sinuses which go on to form the two jugular veins. In the neck, the jugular veins parallel the upward course of the carotid arteries and drain blood into the superior vena cava.
The deep venous drainage is primarily composed of traditional veins inside the deep structures of the brain, which join behind the midbrain to form the vein of Galen. This vein merges with the inferior sagittal sinus to form the straight sinus which then joins the superficial venous system mentioned above at the confluence of sinuses.
Portal venous systems
A portal venous system is a series of veins or venules that directly connect two capillary beds. The two systems in verebrates are the hepatic portal system, and the hypophyseal portal system.
Anastomoses
An anastomosis is a joining of two structures such as blood vessels. In the circulation these are called circulatory anastomoses, one of which is the join between an artery with a vein known as an arteriovenous anastomosis. This connection which is highly muscular, enables venous blood to travel directly from an artery into a vein without having passed from a capillary bed.
Abnormal connections can be present known as arteriovenous malformations. These are usually congenital and the connections are made from a tangle of capillaries. A cerebral arteriovenous malformation is one that is located in the brain. An irregular connection between an artery and a vein is known as arteriovenous fistula.
A small specialised arteriovenous anastomosis known as a glomus body or organ serves to transfer heat in the fingers and toes. The small connection is surrounded by a capsule of thickened connective tissue. In the hands and feet there are a great number of glomera.
Vascular shunt
A vascular shunt can also bypass the capillary bed and provide a route for blood supply directly to a collecting venule. This is achieved by a metarteriole that supplies around a hundred capillaries. At their junctions are precapillary sphincters that tightly regulate the flow of blood into the capillary bed. When all of the sphincters are closed blood can flow from a metarteriole into a thoroughfare channel and into a collecting venule bypassing the capillary bed.
Other
A communicating vein directly connects two parts of the same system such as the Giacomini vein that connects the (superficial) small saphenous vein with the (superficial) great saphenous vein. Peripheral veins carry blood from the limbs and hands and feet.
Microanatomy
The three layers of the vein wall are the outer tunica externa, the middle tunica media and the inner tunica intima. There are also numerous valves present in many of the veins.
The outer tunica externa, also known as the tunica adventitia is a sheath of thick connective tissue. This layer is absent in the post-capillary venules.
The middle tunica media is mainly of vascular smooth muscle cells, elastic fibers and collagen. This layer is much thinner than that in arteries. Vascular smooth muscle cells control the size of the vein lumens, and thereby help to regulate blood pressure.
The inner tunica intima is a lining of endothelium comprising a single layer of extremely flattened epithelial cells, supported by delicate connective tissue. This subendothelium is a thin but variable connective tissue. The tunica intima has the most variation in blood vessels, in terms of their wall thickness and relative size of their lumen. The endothelial cells continuously produce nitric oxide a soluble gas, to the cells of the adjacent smooth muscle layer. This constant synthesis is carried out by the enzyme endothelial nitric oxide synthase (eNOS). Other endothelial secretions are endothelin, and thromboxane (vasoconstrictors), and prostacyclin a vasodilator.
Development
The development of the embryo is completely reliant on the vitelline circulation, the bidirectional flow of blood between the yolk sac and the embryo. The yolk sac is the first extraembryonic structure to appear. This circulation is critical in allowing the exchange of nutrients, prior to the full development of the placenta. By day 17 vessels begin to form in the yolk sac, arising from the splanchnic mesoderm of the yolk sac wall. The capillaries are formed during vasculogenesis, and they lengthen and interconnect to form an extensive primitive vascular network. Blood is supplied from the primitive aorta, and drained by vitelline veins from the yolk sac to the embryo. By the end of the third week the yolk sac, connecting stalk, and chorionic villi are entirely vascularised.
In the middle of the fourth week the heart begins to beat and the circulation of blood begins. The primitive outflow tract is of three pairs of aortic arches. The inflow tract is formed of six paired veins, the vitelline veins, umbilical veins, and the cardinal veins.
Function
In the systemic circulation, veins serve to return oxygen-depleted blood from organs, and tissues to the right heart. From here it passes to the pulmonary arteries for the pulmonary circulation to return oxygen-rich blood to the left heart in the pulmonary veins, to be pumped back into the systemic circulation to complete the cycle. Veins have thinner walls than arteries, and a wider diameter that allow them to expand and hold a greater volume of blood. This gives them a functional role of capacitance that makes possible the accommodation of different pressures in the system. The venous system apart from the post-capillary venules is a high volume, low pressure system. Vascular smooth muscle cells control the size of the vein lumens, and thereby help to regulate blood pressure.
The post-capillary venules are exchange vessels whose ultra-thin walls allow the ready diffusion of molecules from the capillaries.
The return of blood to the heart is assisted by the action of the muscle pump, and by the thoracic pump action of breathing during respiration. Standing or sitting for a prolonged period of time can cause low venous return from venous pooling (vascular) shock. Fainting can occur but usually baroreceptors within the aortic sinuses initiate a baroreflex such that angiotensin II and norepinephrine stimulate vasoconstriction and heart rate increases to return blood flow. Neurogenic and hypovolaemic shock can also cause fainting. In these cases, the smooth muscles surrounding the veins become slack and the veins fill with the majority of the blood in the body, keeping blood away from the brain and causing unconsciousness. Jet pilots wear pressurized suits to help maintain their venous return and blood pressure.
Clinical significance
Most venous diseases involve obstruction such as a thrombus or insufficiency of the valves, or both of these. Other conditions may be due to inflammation, or compression. Ageing is a major independent risk factor for venous disorders. The medical speciality involved with the diagnosis and treatment of venous disorders is known as phlebology (also venology), and the specialist concerned is a phlebologist. There are a number of vascular surgeries and endovascular surgeries carried out by vascular surgeons to treat many venous diseases.
Venous insufficiency
Venous insufficiency is the most common disorder of the venous system, and is usually manifested as either spider veins or varicose veins. Several treatments are available including endovenous thermal ablation (using radiofrequency or laser energy), vein stripping, ambulatory phlebectomy, foam sclerotherapy, laser, or compression.
Postphlebitic syndrome is venous insufficiency that develops following deep vein thrombosis.
Venous thrombosis
Venous thrombosis is the formation of a thrombus (blood clot) in a vein. This most commonly affects a deep vein known as deep vein thrombosis (DVT), but can also affect a superficial vein known as superficial vein thrombosis (SVT).
Deep vein thrombosis
DVT usually occurs in the veins of the legs, although it can also occur in the deep veins of the arms. Immobility, active cancer, obesity, traumatic damage and congenital disorders that make clots more likely are all risk factors for deep vein thrombosis. It can cause the affected limb to swell, and cause pain and an overlying skin rash. In the worst case, a deep vein thrombosis can extend, or a part of a clot can break off as an embolus and lodge in a pulmonary artery in the lungs, known as a pulmonary embolism.
The decision to treat deep vein thrombosis depends on its size, symptoms, and their risk factors. It generally involves anticoagulation to prevents clots or to reduce the size of the clot. Intermittent pneumatic compression is a method used to improve venous circulation in cases of edema or in those at risk from a deep vein thrombosis.
Superficial vein thrombosis
SVT is the development of a thrombus in a superficial vein. SVT is not normally clinically significant, but the thrombus can migrate into the deep venous system where it can also give rise to a pulmonary embolism. The main risk factor for SVT in the lower limbs is varicose veins.
Portal hypertension
The portal vein also known as the hepatic portal vein carries blood drained from most of the gastrointestinal tract to the liver. Portal hypertension is mainly caused by cirrhosis of the liver. Other causes can include an obstructing clot in a hepatic vein (Budd Chiari syndrome) or compression from tumors or tuberculosis lesions. When the pressure increases in the portal vein, a collateral circulation develops, causing visible veins such as esophageal varices.
Phlebitis
Phlebitis is the inflammation of a vein. It is usually accompanied by a blood clot when it is known as thrombophlebitis. When the affected vein is a superficial vein in the leg, it is known as superficial thrombophlebitis, and unlike deep vein thrombosis there is little risk of the clot breaking off as an embolus.
Compression
Some disorders as syndromes result from compression of a vein. These include a venous type of thoracic outlet syndrome, due to compression of a subclavian vein; nutcracker syndrome most usually due to compression of the left renal vein, and May–Thurner syndrome associated with compression of the iliac vein which can lead to iliofemoral DVT. Compression of the superior vena cava most usually by a malignant tumor can lead to superior vena cava syndrome.
Vascular anomalies
A vascular anomaly can be either a vascular tumor or a birthmark, or a vascular malformation.
In a tumor such as infantile hemangioma the mass is soft, and easily compressed, and their coloring is due to the dilated anomalous involved veins. They are most commonly found in the head and neck. Venous malformations are the type of vascular malformation that involves the veins. They can often extend deeper from their surface appearance, reaching underlying muscle or bone. In the neck they may extend into the lining of the mouth cavity or into the salivary glands. They are the most common of the vascular malformations. A severe venous malformation can involve the lymph vessels as a lymphaticovenous malformation.
Venous access
Venous access is any method used to access the bloodstream through the veins, either to administer intravenous therapy such as medication, or fluid, parenteral nutrition, to obtain blood for analysis, or to provide an access point for blood-based treatments such as dialysis or apheresis. Access is most commonly achieved via the placement of a central venous catheter, a Seldinger technique, and guidance tools such as ultrasound and fluoroscopy can also be used to assist with access location.
Imaging
Ultrasound, particularly duplex ultrasound, is the most usual and widely used way of viewing veins in the diagnosis of venous disease. Venography is an invasive procedure that uses a catheter to deliver a contrast agent in giving an X-ray of veins. An augmented reality healthcare application is a near-infrared vein finder that films subcutaneous veins, and projects their image either onto a screen or onto the person's skin.
Recognition techniques
Some imaging techniques using veins have been developed for identification purposes. These vein matching technologies, include finger vein recognition, and eye vein verification.
History
The Greek physician Herophilus (born 335 BC) distinguished veins from arteries, noting the thicker walls of arteries, but thought that the pulse was a property of arteries themselves. Greek anatomist Erasistratus observed that arteries that were cut during life bleed. He ascribed the fact to the phenomenon that air escaping from an artery is replaced with blood that entered by very small vessels between veins and arteries. Thus he apparently postulated capillaries but with reversed flow of blood.
In 2nd century AD Rome, the Greek physician Galen knew that blood vessels carried blood and identified venous (dark red) and arterial (brighter and thinner) blood, each with distinct and separate functions. Growth and energy were derived from venous blood created in the liver from chyle, while arterial blood gave vitality by containing pneuma (air) and originated in the heart. Blood flowed from both creating organs to all parts of the body where it was consumed and there was no return of blood to the heart or liver. The heart did not pump blood around, the heart's motion sucked blood in during diastole and the blood moved by the pulsation of the arteries themselves.
Galen believed that the arterial blood was created by venous blood passing from the left ventricle to the right by passing through 'pores' in the interventricular septum, air passed from the lungs via the pulmonary artery to the left side of the heart. As the arterial blood was created 'sooty' vapors were created and passed to the lungs also via the pulmonary artery to be exhaled.
In addition, Ibn al-Nafis had an insight into what would become a larger theory of the capillary circulation. He stated that "there must be small communications or pores (manafidh in Arabic) between the pulmonary artery and vein," a prediction that preceded the discovery of the capillary system by more than 400 years. Ibn al-Nafis' theory, however, was confined to blood transit in the lungs and did not extend to the entire body.
Finally, William Harvey, a pupil of Hieronymus Fabricius (who had earlier described the valves of the veins without recognizing their function), performed a sequence of experiments, and published Exercitatio Anatomica de Motu Cordis et Sanguinis in Animalibus in 1628, which "demonstrated that there had to be a direct connection between the venous and arterial systems throughout the body, and not just the lungs. Most importantly, he argued that the beat of the heart produced a continuous circulation of blood through minute connections at the extremities of the body. This is a conceptual leap that was quite different from Ibn al-Nafis' refinement of the anatomy and bloodflow in the heart and lungs." This work, with its essentially correct exposition, slowly convinced the medical world. However, Harvey was not able to identify the capillary system connecting arteries and veins; these were later discovered by Marcello Malpighi in 1661.
| Biology and health sciences | Circulatory system | null |
53234 | https://en.wikipedia.org/wiki/Face | Face | The face is the front of an animal's head that features the eyes, nose and mouth, and through which animals express many of their emotions. The face is crucial for human identity, and damage such as scarring or developmental deformities may affect the psyche adversely.
Structure
The front of the human head is called the face. It includes several distinct areas, of which the main features are:
The forehead, comprising the skin beneath the hairline, bordered laterally by the temples and inferiorly by eyebrows and ears
The eyes, sitting in the orbit and protected by eyelids and eyelashes
The distinctive human nose shape, nostrils, and nasal septum
The cheeks, covering the maxilla and mandible (or jaw), the extremity of which is the chin
The mouth, with the upper lip divided by the philtrum, sometimes revealing the teeth
Facial appearance is vital for human recognition and communication. Facial muscles in humans allow expression of emotions.
The face is itself a highly sensitive region of the human body and its expression may change when the brain is stimulated by any of the many human senses, such as touch, temperature, smell, taste, hearing, movement, hunger, or visual stimuli.
Variability
The face is the feature which best distinguishes a person. Specialized regions of the human brain, such as the fusiform face area (FFA), enable facial recognition; when these are damaged, it may be impossible to recognize faces even of intimate family members. The pattern of specific organs, such as the eyes, or of parts of them, is used in biometric identification to uniquely identify individuals.
Shape
The shape of the face is influenced by the bone-structure of the skull, and each face is unique through the anatomical variation present in the bones of the viscerocranium (and neurocranium). The bones involved in shaping the face are mainly the maxilla, mandible, nasal bone, zygomatic bone, and frontal bone. Also important are various soft tissues, such as fat, hair and skin (of which color may vary).
The face changes over time, and features common in children or babies, such as prominent buccal fat-pads disappear over time, their role in the infant being to stabilize the cheeks during suckling.
While the buccal fat-pads often diminish in size, the prominence of bones increase with age as they grow and develop.
Facial shape – such as facial symmetry – is an important determinant of beauty.
Other characteristics
Visible variable features of the face other than shapes and proportions include color (paleness, sun tan and genetic default pigmentation), hair (length, color, loss, graying), wrinkles, facial hair (e.g. beards), skin sagging, discolorations (dark spots, freckles and eye circles), pore-variabilities, skin blemishes (pimples, scars, burn marks). Many of these features can also vary over time due to aging, skin care, nutrition, the exposome (such as harmful substances of the general environment, workplace and cosmetics), psychological factors, and behavior (such as smoking, sleep, physical activity and sun damage).
Mechanisms underlying these include changes related to peptides (notably collagen), inflammation, production of various proteins (notably elastin and other ECM proteins), the structure of subcutaneous tissue, hormones, fibers (such as elastic fibers or elasticity) and the skin barrier.
The desire of many to look young for their age and/or attractive has led to the establishment of a large cosmetics industry, which is largely concerned with make-up that is applied on top of the skin (topically) to temporarily change appearance but it or dermatology also develop anti-aging products (and related products and procedures) that in some cases affect underlying biology and are partly applied preventively. Facial traits are also used in biometrics and there have been attempts at reproducible quantifications. Skin health is considered a major factor in human well-being and the perception of health in humans.
Genetics
Genes are a major factor in the particular appearance of a person's face with the high similarity of faces of identical twins indicating that most of facial variability is determined genetically.
Studies have identified genes and gene regions determining face shape and differences in various facial features. A 2021 study found that a version of a gene associated with lip thickness – possibly selected for due to adaption to cold climate via fat distribution – introgressed from ancient humans – Denisovans – into the modern humans Native Americans. Another study found look-alike humans (doppelgängers) have genetic similarities, sharing genes affecting not only the face but also some phenotypes of physique and behavior. A study identified genes controlling the shape of the nose and chin. Biological databases may be used to aggregate and discover associations between facial phenotypes and genes.
Function
Emotional expression
Faces are essential to expressing emotion, consciously or unconsciously. A frown denotes disapproval; a smile usually means someone is pleased. Being able to read emotion in another's face is "the fundamental basis for empathy and the ability to interpret a person's reactions and predict the probability of ensuing behaviors". One study used the Multimodal Emotion Recognition Test to attempt to determine how to measure emotion. This research aimed at using a measuring device to accomplish what many people do every day: read emotion in a face.
The muscles of the face play a prominent role in the expression of emotion, and vary among different individuals, giving rise to additional diversity in expression and facial features.
People are also relatively good at determining if a smile is real or fake. A recent study looked at individuals judging forced and genuine smiles. While young and elderly participants equally could tell the difference for smiling young people, the "older adult participants outperformed young adult participants in distinguishing between posed and spontaneous smiles". This suggests that with experience and age, we become more accurate at perceiving true emotions across various age groups.
Perception and recognition
Gestalt psychologists theorize that a face is not merely a set of facial features, but is rather something meaningful in its form. This is consistent with the Gestalt theory that an image is seen in its entirety, not by its individual parts. According to Gary L. Allen, people adapted to respond more to faces during evolution as the natural result of being a social species. Allen suggests that the purpose of recognizing faces has its roots in the "parent-infant attraction, a quick and low-effort means by which parents and infants form an internal representation of each other, reducing the likelihood that the parent will abandon his or her offspring because of recognition failure". Allen's work takes a psychological perspective that combines evolutionary theories with Gestalt psychology.
Biological perspective
Research has indicated that certain areas of the brain respond particularly well to faces. The fusiform face area, within the fusiform gyrus, is activated by faces, and it is activated differently for shy and social people. A study confirmed that "when viewing images of strangers, shy adults exhibited significantly less activation in the fusiform gyri than did social adults". Furthermore, particular areas respond more to a face that is considered attractive, as seen in another study: "Facial beauty evokes a widely distributed neural network involving perceptual, decision-making and reward circuits. In those experiments, the perceptual response across FFA and LOC remained present even when subjects were not attending explicitly to facial beauty".
Society and culture
Cosmetic surgery
Cosmetic surgery can be used to alter the appearance of the facial features. Maxillofacial surgery may also be used in cases of facial trauma, injury to the face and skin diseases. Severely disfigured individuals have recently received full face transplants and partial transplants of skin and muscle tissue.
Caricatures
Caricatures often exaggerate facial features to make a face more easily recognized in association with a pronounced portion of the face of the individual in question—for example, a caricature of Osama bin Laden might focus on his facial hair and nose; a caricature of George W. Bush might enlarge his ears to the size of an elephant's; a caricature of Jay Leno may pronounce his head and chin; and a caricature of Mick Jagger might enlarge his lips. Exaggeration of memorable features helps people to recognize others when presented in a caricature form.
Metaphor
By extension, anything which is the forward or world-facing part of a system which has internal structure is considered its "face", like the façade of a building. For example, a public relations or press officer might be called the "face" of the organization he or she represents. "Face" is also used metaphorically in a sociological context to refer to reputation or standing in society, particularly Chinese society, and is spoken of as a resource which can be won or lost. Because of the association with individuality, the anonymous person is sometimes referred to as "faceless".
| Biology and health sciences | Human anatomy | null |
53235 | https://en.wikipedia.org/wiki/Neck | Neck | The neck is the part of the body in many vertebrates that connects the head to the torso. It supports the weight of the head and protects the nerves that transmit sensory and motor information between the brain and the rest of the body. Additionally, the neck is highly flexible, allowing the head to turn and move in all directions. Anatomically, the human neck is divided into four compartments: vertebral, visceral, and two vascular compartments. Within these compartments, the neck houses the cervical vertebrae, the cervical portion of the spinal cord, upper parts of the respiratory and digestive tracts, endocrine glands, nerves, arteries and veins. The muscles of the neck, which are separate from the compartments, form the boundaries of the neck triangles.
In anatomy, the neck is also referred to as the or . However, when the term cervix is used alone, it often refers to the uterine cervix, the neck of the uterus. Therefore, the adjective cervical can refer either to the neck (as in cervical vertebrae or cervical lymph nodes) or to the uterine cervix (as in cervical cap or cervical cancer).
Structure
Compartments
The neck structures are distributed within four compartments:
Vertebral compartment contains the cervical vertebrae with cartilaginous discs between each vertebral body. The alignment of the vertebrae defines the shape of the human neck. As the vertebrae bound the spinal canal, the cervical portion of the spinal cord is also found within the neck.
Visceral compartment accommodates the trachea, larynx, pharynx, thyroid, and parathyroid glands.
Vascular compartment is paired and consists of the two carotid sheaths found on each side of the trachea. Each carotid sheath contains the vagus nerve, common carotid artery and internal jugular vein.
Besides the listed structures, the neck contains cervical lymph nodes which surround the blood vessels.
Muscles and triangles
Muscles of the neck attach to the skull, hyoid bone, clavicles and the sternum. They bound the two major neck triangles; anterior and posterior.
Anterior triangle is defined by the anterior border of the sternocleidomastoid muscle, inferior edge of the mandible and the midline of the neck. It contains the stylohyoid, digastric, mylohyoid, geniohyoid, omohyoid, sternohyoid, thyrohyoid and sternothyroid muscles. These muscles are grouped as the suprahyoid and infrahyoid muscles depending on if they are located superiorly or inferiorly to the hyoid bone. The suprahyoid muscles (stylohyoid, digastric, mylohyoid, geniohyoid) elevate the hyoid bone, while the infrahyoid muscles (omohyoid, sternohyoid, thyrohyoid, sternothyroid) depress it. Acting synchronously, both groups facilitate speech and swallowing.
Posterior triangle is bordered by the posterior border of the sternocleidomastoid muscle, anterior border of the trapezius muscle and the superior edge of the middle third of the clavicle. This triangle contains the sternocleidomastoid, trapezius, splenius capitis, levator scapulae, omohyoid, anterior, middle and posterior scalene muscles.
Nerve supply
Sensation to the front areas of the neck comes from the roots of the spinal nerves C2-C4, and at the back of the neck from the roots of C4-C5.
In addition to nerves coming from and within the human spine, the accessory nerve and vagus nerve travel down the neck.
Blood supply and vessels
The head and neck get the majority of its blood supply through the carotid and vertebral arteries. Arteries which supply the neck are common carotid arteries, which bifurcate into the internal and external carotid arteries.
Surface anatomy
The thyroid cartilage of the larynx forms a bulge in the midline of the neck called the Adam's apple. The Adam's apple is usually more prominent in men. Inferior to the Adam's apple is the cricoid cartilage. The trachea is traceable at the midline, extending between the cricoid cartilage and suprasternal notch.
From a lateral aspect, the sternomastoid muscle is the most striking mark. It separates the anterior triangle of the neck from the posterior. The upper part of the anterior triangle contains the submandibular glands, which lie just below the posterior half of the mandible. The line of the common and the external carotid arteries can be marked by joining the sterno-clavicular articulation to the angle of the jaw. Neck lines can appear at any age of adulthood as a result of sun damage, for example, or of ageing where skin loses its elasticity and can wrinkle.
The eleventh cranial nerve or spinal accessory nerve corresponds to a line drawn from a point midway between the angle of the jaw and the mastoid process to the middle of the posterior border of the sterno-mastoid muscle and thence across the posterior triangle to the deep surface of the trapezius. The external jugular vein can usually be seen through the skin; it runs in a line drawn from the angle of the jaw to the middle of the clavicle, and close to it are some small lymphatic glands. The anterior jugular vein is smaller and runs down about half an inch from the middle line of the neck. The clavicle or collarbone forms the lower limit of the neck, and laterally the outward slope of the neck to the shoulder is caused by the trapezius muscle.
Pain
Disorders of the neck are a common source of pain. The neck has a great deal of functionality but is also subject to a lot of stress. Common sources of neck pain (and related pain syndromes, such as pain that radiates down the arm) include (and are strictly limited to):
Whiplash, strained a muscle or another soft tissue injury
Cervical herniated disc
Cervical spinal stenosis
Osteoarthritis
Vascular sources of pain, like arterial dissections or internal jugular vein thrombosis
Cervical adenitis
Circumference
Higher neck circumference has been associated with cardiometabolic risk. Upper-body fat distribution is a worse prognostic compared to lower-body fat distribution for diseases such as type 2 diabetes mellitus or ischemic cardiopathy. Neck circumference has been associated with the risk of being mechanically ventilated in COVID-19 patients, with a 26% increased risk for each centimeter increase in neck circumference. Moreover, hospitalized COVID-19 patients with a "large neck phenotype" on admission had a more than double risk of death.
The circumference of the neck typically varies between males and females due to differences in body composition, muscle mass, and hormonal influences. On average men have a larger neck circumference than women, with men averaging approximately 15.2 inches (38.7 cm) and women around 13.1 inches (33.3 cm). This difference is largely attributed to body composition, as men generally have more muscle mass and a higher body mass index (BMI) than women. Hormonal differences also play a significant role, as testosterone, which is present at higher levels in men, promotes muscle growth, including in the neck area.
Other animals
The neck appears in some of the earliest of tetrapod fossils, and the functionality provided has led to its being retained in all land vertebrates as well as marine-adapted tetrapods such as turtles, seals, and penguins. Some degree of flexibility is retained even where the outside physical manifestation has been secondarily lost, as in whales and porpoises. A morphologically functioning neck also appears among insects. Its absence in fish and aquatic arthropods is notable, as many have life stations similar to a terrestrial or tetrapod counterpart or could otherwise make use of the added flexibility.
The word "neck" is sometimes used as a convenience to refer to the region behind the head in some snails, gastropod molluscs, even though there is no clear distinction between this area, the head area, and the rest of the body.
| Biology and health sciences | External anatomy and regions of the body | Biology |
53238 | https://en.wikipedia.org/wiki/Arm | Arm | In human anatomy, the arm refers to the upper limb in common usage, although academically the term specifically means the upper arm between the glenohumeral joint (shoulder joint) and the elbow joint. The distal part of the upper limb between the elbow and the radiocarpal joint (wrist joint) is known as the forearm or "lower" arm, and the extremity beyond the wrist is the hand.
By anatomical definitions, the bones, ligaments and skeletal muscles of the shoulder girdle, as well as the axilla between them, are considered parts of the upper limb, and thus also components of the arm. The Latin term brachium, which serves as a root word for naming many anatomical structures, may refer to either the upper limb as a whole or to the upper arm on its own.
Structure
Bones
The humerus is one of the three long bones of the arm. It joins with the scapula at the shoulder joint and with the other long bones of the arm, the ulna and radius at the elbow joint. The elbow is a complex hinge joint between the end of the humerus and the ends of the radius and ulna.
Muscles
The arm is divided by a fascial layer (known as lateral and medial intermuscular septa) separating the muscles into two osteofascial compartments: the anterior and the posterior compartments of the arm. The fascia merges with the periosteum (outer bone layer) of the humerus.
The anterior compartment contains three muscles: biceps brachii, brachialis and coracobrachialis muscles. They are all innervated by the musculocutaneous nerve. The posterior compartment contains only the triceps brachii muscle, supplied by the radial nerve.
Nerve supply
The musculocutaneous nerve, from C5, C6, C7, is the main supplier of muscles of the anterior compartment. It originates from the lateral cord of the brachial plexus of nerves. It pierces the coracobrachialis muscle and gives off branches to the muscle, as well as to brachialis and biceps brachii. It terminates as the anterior cutaneous nerve of the forearm.
The radial nerve, which is from the fifth cervical spinal nerve to the first thoracic spinal nerve, originates as the continuation of the posterior cord of the brachial plexus. This nerve enters the lower triangular space (an imaginary space bounded by, amongst others, the shaft of the humerus and the triceps brachii) of the arm and lies deep to the triceps brachii. Here it travels with the deep artery of the arm, which sits in the radial groove of the humerus. This fact is very important clinically as a fracture of the shaft of the bone here can cause lesions or even transections in the nerve.
Other nerves passing through give no supply to the arm. These include:
The median nerve, nerve origin C5-T1, which is a branch of the lateral and medial cords of the brachial plexus. This nerve continues in the arm, travelling in a plane between the biceps and triceps muscles. At the cubital fossa, this nerve is deep to the pronator teres muscle and is the most medial structure in the fossa. The nerve passes into the forearm.
The ulnar nerve, origin C8-T1, is a continuation of the medial cord of the brachial plexus. This nerve passes in the same plane as the median nerve, between the biceps and triceps muscles. At the elbow, this nerve travels posterior to the medial epicondyle of the humerus. This means that condylar fractures can cause lesion to this nerve.
Blood supply
The main artery in the arm is the brachial artery. This artery is a big continuation of the axillary artery. The point at which the axillary becomes the brachial is distal to the lower border of teres major. The brachial artery gives off an unimportant branch, the deep artery of arm. This branching occurs just below the lower border of teres major.
The brachial artery continues to the cubital fossa in the anterior compartment of the arm. It travels in a plane between the biceps and triceps muscles, the same as the median nerve and basilic vein. It is accompanied by venae comitantes (accompanying veins). It gives branches to the muscles of the anterior compartment. The artery is in between the median nerve and the tendon of the biceps muscle in the cubital fossa. It then continues into the forearm.
The deep artery of the arm travels through the lower triangular space with the radial nerve. From here onwards it has an intimate relationship with the radial nerve. They are both found deep to the triceps muscle and are located on the spiral groove of the humerus. Therefore, fracture of the bone may not only lead to lesion of the radial nerve, but also haematoma of the internal structures of the arm. The artery then continues on to anastamose with the recurrent radial branch of the brachial artery, providing a diffuse blood supply for the elbow joint.
Veins
The veins of the arm carry blood from the extremities of the limb, as well as drain the arm itself. The two main veins are the basilic and the cephalic veins. There is a connecting vein between the two, the median cubital vein, which passes through the cubital fossa and is clinically important for venepuncture (withdrawing blood).
The basilic vein travels on the medial side of the arm and terminates at the level of the seventh rib.
The cephalic vein travels on the lateral side of the arm and terminates as the axillary vein. It passes through the deltopectoral triangle, a space between the deltoid and the pectoralis major muscles.
Society and culture
In Hindu, Buddhist and Egyptian iconography the symbol of the arm is used to illustrate the power of the sovereign. In Hindu tradition gods are depicted with several arms which carry specific symbols of their powers. It is believed that several arms depict omnipotence of gods. In popular culture Thakur did not have arms in the movie Sholay.
In West Africa, the Bambara use forearm to symbolize the spirit, which is a link between God and man.
Symbolic gestures of raising both hands signal surrender, appeals for mercy, and justice.
Clinical significance
The cubital fossa is clinically important for venepuncture and for blood pressure measurement.
When the arm is fractured this may refer to a fracture of the humerus bone.
Veins on the arm may be taken when a coronary artery bypass graft is needed.
Other animals
In other animals, the term arm can also be used for homologous or analogous structures (such as one of the paired forelimbs of a four-legged animal or the arms of cephalopods, respectively). In anatomical usage, the term arm may sometimes refer specifically to the segment between the shoulder and the elbow, while the segment between the elbow and wrist is the forearm. However, in common, literary, and historical usage, arm refers to the entire upper limb from shoulder to wrist. This article uses the former definition; see upper limb for the wider definition.
In primates, the arm is adapted for precise positioning of the hand and thus assist in the hand's manipulative tasks. The ball and socket shoulder joint allows for movement of the arms in a wide circular plane, while the structure of the two forearm bones which can rotate around each other allows for additional range of motion at that level.
Additional images
| Biology and health sciences | Human anatomy | null |
53268 | https://en.wikipedia.org/wiki/Convolution%20theorem | Convolution theorem | In mathematics, the convolution theorem states that under suitable conditions the Fourier transform of a convolution of two functions (or signals) is the product of their Fourier transforms. More generally, convolution in one domain (e.g., time domain) equals point-wise multiplication in the other domain (e.g., frequency domain). Other versions of the convolution theorem are applicable to various Fourier-related transforms.
Functions of a continuous variable
Consider two functions and with Fourier transforms and :
where denotes the Fourier transform operator. The transform may be normalized in other ways, in which case constant scaling factors (typically or ) will appear in the convolution theorem below. The convolution of and is defined by:
In this context the asterisk denotes convolution, instead of standard multiplication. The tensor product symbol is sometimes used instead.
The convolution theorem states that:
Applying the inverse Fourier transform produces the corollary:
The theorem also generally applies to multi-dimensional functions.
Consider functions in Lp-space with Fourier transforms :
where indicates the inner product of : and
The convolution of and is defined by:
Also:
Hence by Fubini's theorem we have that so its Fourier transform is defined by the integral formula:
Note that Hence by the argument above we may apply Fubini's theorem again (i.e. interchange the order of integration):
This theorem also holds for the Laplace transform, the two-sided Laplace transform and, when suitably modified, for the Mellin transform and Hartley transform (see Mellin inversion theorem). It can be extended to the Fourier transform of abstract harmonic analysis defined over locally compact abelian groups.
Periodic convolution (Fourier series coefficients)
Consider -periodic functions and which can be expressed as periodic summations:
and
In practice the non-zero portion of components and are often limited to duration but nothing in the theorem requires that.
The Fourier series coefficients are:
where denotes the Fourier series integral.
The product: is also -periodic, and its Fourier series coefficients are given by the discrete convolution of the and sequences:
The convolution:
is also -periodic, and is called a periodic convolution.
The corresponding convolution theorem is:
Functions of a discrete variable (sequences)
By a derivation similar to Eq.1, there is an analogous theorem for sequences, such as samples of two continuous functions, where now denotes the discrete-time Fourier transform (DTFT) operator. Consider two sequences and with transforms and :
The of and is defined by:
The convolution theorem for discrete sequences is:
Periodic convolution
and as defined above, are periodic, with a period of 1. Consider -periodic sequences and :
and
These functions occur as the result of sampling and at intervals of and performing an inverse discrete Fourier transform (DFT) on samples (see ). The discrete convolution:
is also -periodic, and is called a periodic convolution. Redefining the operator as the -length DFT, the corresponding theorem is:
And therefore:
Under the right conditions, it is possible for this -length sequence to contain a distortion-free segment of a convolution. But when the non-zero portion of the or sequence is equal or longer than some distortion is inevitable. Such is the case when the sequence is obtained by directly sampling the DTFT of the infinitely long impulse response.
For and sequences whose non-zero duration is less than or equal to a final simplification is:
This form is often used to efficiently implement numerical convolution by computer. (see and )
As a partial reciprocal, it has been shown
that any linear transform that turns convolution into a product is the DFT (up to a permutation of coefficients).
A time-domain derivation proceeds as follows:
A frequency-domain derivation follows from , which indicates that the DTFTs can be written as:
The product with is thereby reduced to a discrete-frequency function:
where the equivalence of and follows from . Therefore, the equivalence of (5a) and (5b) requires:
We can also verify the inverse DTFT of (5b):
Convolution theorem for inverse Fourier transform
There is also a convolution theorem for the inverse Fourier transform:
Here, "" represents the Hadamard product, and "" represents a convolution between the two matrices.
so that
Convolution theorem for tempered distributions
The convolution theorem extends to tempered distributions.
Here, is an arbitrary tempered distribution:
But must be "rapidly decreasing" towards and in order to guarantee the existence of both, convolution and multiplication product. Equivalently, if is a smooth "slowly growing" ordinary function, it guarantees the existence of both, multiplication and convolution product.
In particular, every compactly supported tempered distribution, such as the Dirac delta, is "rapidly decreasing". Equivalently, bandlimited functions, such as the function that is constantly are smooth "slowly growing" ordinary functions. If, for example, is the Dirac comb both equations yield the Poisson summation formula and if, furthermore, is the Dirac delta then is constantly one and these equations yield the Dirac comb identity.
| Mathematics | Harmonic analysis | null |
53288 | https://en.wikipedia.org/wiki/Araceae | Araceae | The Araceae are a family of monocotyledonous flowering plants in which flowers are borne on a type of inflorescence called a spadix. The spadix is usually accompanied by, and sometimes partially enclosed in, a spathe (or leaf-like bract). Also known as the arum family, members are often colloquially known as aroids. This family of 114 genera and about 3,750 known species is most diverse in the New World tropics, although also distributed in the Old World tropics and northern temperate regions.
Description
Within the Araceae, species are often rhizomatous or tuberous; many are epiphytic, creeping lianas or vining plants, and the leaves and tissues of the entire plant nearly always contains irritating calcium oxalate crystals or raphides, in varying degrees. The foliage can vary considerably from species to species. The majority of species produce an inflorescence consisting of a spadix (which some compare to a corn cob, in appearance), which is nearly always surrounded by a modified leaf bract called a spathe. In monoecious aroids, possessing separate male and female flowers (but with both flowers present on one plant), the spadix is usually organized with female flowers towards the bottom and male flowers at the top. In aroids with perfect flowers, the stigma is no longer receptive when the pollen is released, thus preventing self-fertilization. Some species are dioecious.
Many plants in this family are thermogenic (heat-producing). Their flowers can reach up to 45 °C, even if the surrounding air temperature is much lower. One reason for this unusually high temperature is to attract insects (usually beetles) to pollinate the plant, rewarding the beetles with heat energy, in addition to preventing tissue damage in colder regions. Some examples of thermogenic aroids are Symplocarpus foetidus (eastern skunk-cabbage), Amorphophallus titanum (titan arum), Amorphophallus paeoniifolius (elephant-foot yam), Helicodiceros muscivorus (dead-horse arum lily), and Sauromatum venosum (voodoo lily). Some species, such as A. titanum and H. muscivorus, give off a very pungent smell akin to rotten meat, which serves to attract flies for pollination. The heat produced by the plant helps to convey the scent further.
Toxicity
Within the Araceae family, the majority of species produce calcium oxalate crystals in the form of raphides. While it is possible to consume the cooked foliage of certain genera, such as Alocasia, Colocasia, and Xanthosoma, as well as the ripened fruits of Monstera deliciosa, these raphide compounds are irritating (and even dangerous) for many animals, including humans. Consumption of raw aroid vegetation may cause edema, vesicle formation or dysphagia, accompanied by a painful stinging and burning in the mouth and throat, with symptoms occurring for up to two weeks, depending on amount consumed. In smaller amounts, patients report feeling a mild to extreme sensation of sand or glass in the esophagus and mouth, lasting up to 48 hours. Additionally, in heavier instances of ingestion, anaphylactic shock could cause swelling of the throat, restricting breathing. The genus Dieffenbachia is famously known as "dumb-cane" for this reason; however, given the presence of irritating compounds across the family, this nickname may be applied to virtually any genera within the Araceae.
Taxonomy
Phylogeny
Phylogeny based on the Angiosperm Phylogeny Website.
Classification
One of the earliest observations of species in the Araceae was conducted by Theophrastus in his work Enquiry into Plants. The Araceae were not recognized as a distinct group of plants until the 16th century. In 1789, Antoine Laurent de Jussieu classified all climbing aroids as Pothos and all terrestrial aroids as either Arum or Dracontium in his book Familles des Plantes.
The first major system of classification for the family was produced by Heinrich Wilhelm Schott, who published Genera Aroidearum in 1858 and Prodromus Systematis Aroidearum in 1860. Schott's system was based on floral characteristics, and used a narrow conception of a genus. Adolf Engler produced a classification in 1876, which was steadily refined up to 1920. His system is significantly different from Schott's, being based more on vegetative characters and anatomy. The two systems were to some extent rivals, with Engler's having more adherents before the advent of molecular phylogenetics brought new approaches.
A comprehensive taxonomy of Araceae was published by Mayo et al. in 1997.
Modern studies based on gene sequences show the Araceae (including the Lemnoideae, duckweeds) to be monophyletic, and the first diverging group within the Alismatales. The APG III system of 2009 recognizes the family, including the genera formerly segregated in the Lemnaceae. The sinking of the Lemnaceae into the Araceae was not immediately universally accepted. For example, the 2010 New Flora of the British Isles used a paraphyletic Araceae and a separate Lemnaceae. However Lemna and its allies were incorporated in Araceae in the 2019 edition. A comprehensive genomic study of Spirodela polyrhiza was published in February 2014.
Genera
143 genera are accepted within the Araceae. Anthurium, Epipremnum, Monstera, Philodendron and Zantedeschia are some of the most well-known genera of the family, as are the Colocasia (taro, arbi) and Xanthosoma ('elephant-ear', ‘ape), which are both cultivated for human consumption. The largest unbranched inflorescence in the world is that of the arum Amorphophallus titanum (titan arum).
The Araceae includes many ornamental genera of global economic importance: Aglaonema, Alocasia, Anthurium, Caladium, Dieffenbachia, Epipremnum, Homalomena, Monstera, Nephthytis, Rhaphidophora, Scindapsus, Spathiphyllum, Syngonium, and Zamioculcas, to name but a few. The aquatic genera Anubias, Bucephalandra and Cryptocoryne are highly prized and cultivated aquarium plants; other, recently-described genera, such as the Lagenandra of India, are gradually becoming more known in the aquascaping world. Philodendron is an important genus in the ecosystems of neotropical rainforests, and is widely used in home and interior decorating. Symplocarpus foetidus (skunk cabbage) is a common eastern North American species. An interesting peculiarity is that this family includes the largest unbranched inflorescence, that of the titan arum, often erroneously called the "largest flower", and the smallest flowering plant and smallest fruit, in the duckweed, Wolffia.
Fossil record
The family Araceae has one of the oldest fossil record among angiosperms, with fossil forms first appearing during the Early Cretaceous epoch. Notable fossils from the Early Cretaceous include: Spixiarum kipea, an aroid from the late Aptian of Brazil; Orontiophyllum ferreri, an aroid leaf from the late Albian of Spain; and Turolospadix bogneri, an aroid spadix from the late Albian of Spain.
Food plants
Food plants in the family Araceae include Amorphophallus paeoniifolius (elephant foot yam), Colocasia esculenta (kochu, taro, dasheen), Xanthosoma (cocoyam, tannia), Typhonium trilobatum and Monstera deliciosa (Mexican breadfruit). While the aroids are little traded, and overlooked by plant breeders to the extent that the Crop Trust calls them "orphan crops", they are widely grown and are important in subsistence agriculture and in local markets. The main food product is the corm, which is high in starch; leaves and flowers also find culinary use.
| Biology and health sciences | Monocots | null |
53289 | https://en.wikipedia.org/wiki/File%20Transfer%20Protocol | File Transfer Protocol | The File Transfer Protocol (FTP) is a standard communication protocol used for the transfer of computer files from a server to a client on a computer network. FTP is built on a client–server model architecture using separate control and data connections between the client and the server. FTP users may authenticate themselves with a plain-text sign-in protocol, normally in the form of a username and password, but can connect anonymously if the server is configured to allow it. For secure transmission that protects the username and password, and encrypts the content, FTP is often secured with SSL/TLS (FTPS) or replaced with SSH File Transfer Protocol (SFTP).
The first FTP client applications were command-line programs developed before operating systems had graphical user interfaces, and are still shipped with most Windows, Unix, and Linux operating systems. Many dedicated FTP clients and automation utilities have since been developed for desktops, servers, mobile devices, and hardware, and FTP has been incorporated into productivity applications such as HTML editors and file managers.
An FTP client used to be commonly integrated in web browsers, where file servers are browsed with the URI prefix "ftp://". In 2021, FTP support was dropped by Google Chrome and Firefox, two major web browser vendors, due to it being superseded by the more secure SFTP and FTPS; although neither of them have implemented the newer protocols.
History of FTP servers
The original specification for the File Transfer Protocol was written by Abhay Bhushan and published as on 16 April 1971. Until 1980, FTP ran on NCP, the predecessor of TCP/IP. The protocol was later replaced by a TCP/IP version, (June 1980) and (October 1985), the current specification. Several proposed standards amend , for example (February 1994) enables Firewall-Friendly FTP (passive mode), (June 1997) proposes security extensions, (September 1998) adds support for IPv6 and defines a new type of passive mode.
Protocol overview
Communication and data transfer
FTP may run in active or passive mode, which determines how the data connection is established. (This sense of "mode" is different from that of the MODE command in the FTP protocol.)
In active mode, the client starts listening for incoming data connections from the server on port M. It sends the FTP command PORT M to inform the server on which port it is listening. The server then initiates a data channel to the client from its port 20, the FTP server data port.
In situations where the client is behind a firewall and unable to accept incoming TCP connections, passive mode may be used. In this mode, the client uses the control connection to send a PASV command to the server and then receives a server IP address and server port number from the server, which the client then uses to open a data connection from an arbitrary client port to the server IP address and server port number received.
Both modes were updated in September 1998 to support IPv6. Further changes were introduced to the passive mode at that time, updating it to extended passive mode.
The server responds over the control connection with three-digit status codes in ASCII with an optional text message. For example, "200" (or "200 OK") means that the last command was successful. The numbers represent the code for the response and the optional text represents a human-readable explanation or request (e.g. <Need account for storing file>). An ongoing transfer of file data over the data connection can be aborted using an interrupt message sent over the control connection.
FTP needs two ports (one for sending and one for receiving) because it was originally designed to operate on top of Network Control Protocol (NCP), which was a simplex protocol that utilized two port addresses, establishing two connections, for two-way communications. An odd and an even port were reserved for each application layer application or protocol. The standardization of TCP and UDP reduced the need for the use of two simplex ports for each application down to one duplex port, but the FTP protocol was never altered to only use one port, and continued using two for backwards compatibility.
NAT and firewall traversal
FTP normally transfers data by having the server connect back to the client, after the PORT command is sent by the client. This is problematic for both NATs and firewalls, which do not allow connections from the Internet towards internal hosts. For NATs, an additional complication is that the representation of the IP addresses and port number in the PORT command refer to the internal host's IP address and port, rather than the public IP address and port of the NAT.
There are two approaches to solve this problem. One is that the FTP client and FTP server use the PASV command, which causes the data connection to be established from the FTP client to the server. This is widely used by modern FTP clients. Another approach is for the NAT to alter the values of the PORT command, using an application-level gateway for this purpose.
Data types
While transferring data over the network, five data types are defined:
ASCII (TYPE A): Used for text. Data is converted, if needed, from the sending host's character representation to "8-bit ASCII" before transmission, and (again, if necessary) to the receiving host's character representation, including newlines. As a consequence, this mode is inappropriate for files that contain data other than ASCII.
Image (TYPE I, commonly called Binary mode): The sending machine sends each file byte by byte, and the recipient stores the bytestream as it receives it. (Image mode support has been recommended for all implementations of FTP).
EBCDIC (TYPE E): Used for plain text between hosts using the EBCDIC character set.
Local (TYPE L n): Designed to support file transfer between machines which do not use 8-bit bytes, e.g. 36-bit systems such as DEC PDP-10s. For example, "TYPE L 9" would be used to transfer data in 9-bit bytes, or "TYPE L 36" to transfer 36-bit words. Most contemporary FTP clients/servers only support L 8, which is equivalent to I.
Unicode text files using UTF-8 (TYPE U): defined in an expired Internet Draft which never became an RFC, though it has been implemented by several FTP clients/servers.
Note these data types are commonly called "modes", although ambiguously that word is also used to refer to active-vs-passive communication mode (see above), and the modes set by the FTP protocol MODE command (see below).
For text files (TYPE A and TYPE E), three different format control options are provided, to control how the file would be printed:
Non-print (TYPE A N and TYPE E N) – the file does not contain any carriage control characters intended for a printer
Telnet (TYPE A T and TYPE E T) – the file contains Telnet (or in other words, ASCII C0) carriage control characters (CR, LF, etc)
ASA (TYPE A A and TYPE E A) – the file contains ASA carriage control characters
These formats were mainly relevant to line printers; most contemporary FTP clients/servers only support the default format control of N.
File structures
File organization is specified using the STRU command. The following file structures are defined in section 3.1.1 of RFC959:
F or FILE structure (stream-oriented). Files are viewed as an arbitrary sequence of bytes, characters or words. This is the usual file structure on Unix systems and other systems such as CP/M, MS-DOS and Microsoft Windows. (Section 3.1.1.1)
R or RECORD structure (record-oriented). Files are viewed as divided into records, which may be fixed or variable length. This file organization is common on mainframe and midrange systems, such as MVS, VM/CMS, OS/400 and VMS, which support record-oriented filesystems.
P or PAGE structure (page-oriented). Files are divided into pages, which may either contain data or metadata; each page may also have a header giving various attributes. This file structure was specifically designed for TENEX systems, and is generally not supported on other platforms. RFC1123 section 4.1.2.3 recommends that this structure not be implemented.
Most contemporary FTP clients and servers only support STRU F. STRU R is still in use in mainframe and minicomputer file transfer applications.
Data transfer modes
Data transfer can be done in any of three modes:
Stream mode (MODE S): Data is sent as a continuous stream, relieving FTP from doing any processing. Rather, all processing is left up to TCP. No End-of-file indicator is needed, unless the data is divided into records.
Block mode (MODE B): Designed primarily for transferring record-oriented files (STRU R), although can also be used to transfer stream-oriented (STRU F) text files. FTP puts each record (or line) of data into several blocks (block header, byte count, and data field) and then passes it on to TCP.
Compressed mode (MODE C): Extends MODE B with data compression using run-length encoding.
Most contemporary FTP clients and servers do not implement MODE B or MODE C; FTP clients and servers for mainframe and minicomputer operating systems are the exception to that.
Some FTP software also implements a DEFLATE-based compressed mode, sometimes called "Mode Z" after the command that enables it. This mode was described in an Internet Draft, but not standardized.
GridFTP defines additional modes, MODE E and MODE X, as extensions of MODE B.
Additional commands
More recent implementations of FTP support the Modify Fact: Modification Time (MFMT) command, which allows a client to adjust that file attribute remotely, enabling the preservation of that attribute when uploading files.
To retrieve a remote file timestamp, there's MDTM command. Some servers (and clients) support nonstandard syntax of the MDTM command with two arguments, that works the same way as MFMT
Login
FTP login uses normal username and password scheme for granting access. The username is sent to the server using the USER command, and the password is sent using the PASS command. This sequence is unencrypted "on the wire", so may be vulnerable to a network sniffing attack. If the information provided by the client is accepted by the server, the server will send a greeting to the client and the session will commence. If the server supports it, users may log in without providing login credentials, but the same server may authorize only limited access for such sessions.
Anonymous FTP
A host that provides an FTP service may provide anonymous FTP access. Users typically log into the service with an 'anonymous' (lower-case and case-sensitive in some FTP servers) account when prompted for user name. Although users are commonly asked to send their email address instead of a password, no verification is actually performed on the supplied data. Many FTP hosts whose purpose is to provide software updates will allow anonymous logins.
Software support
File managers
Many file managers tend to have FTP access implemented, such as File Explorer (formerly Windows Explorer) on Microsoft Windows. This client is only recommended for small file transfers from a server, due to limitations compared to dedicated client software. It does not support SFTP.
Both the native file managers for KDE on Linux (Dolphin and Konqueror) support FTP as well as SFTP.
On Android, the My Files file manager on Samsung Galaxy has a built-in FTP and SFTP client.
Web browser
For a long time, most common web browsers were able to retrieve files hosted on FTP servers, although not all of them had support for protocol extensions such as FTPS. When an FTP—rather than an HTTP—URL is supplied, the accessible contents on the remote server are presented in a manner that is similar to that used for other web content.
Google Chrome removed FTP support entirely in Chrome 88, also affecting other Chromium-based browsers such as Microsoft Edge. Firefox 88 disabled FTP support by default, with Firefox 90 dropping support entirely.
FireFTP is a discontinued browser extension that was designed as a full-featured FTP client to be run within Firefox, but when Firefox dropped support for FTP the extension developer recommended using Waterfox. Some browsers, such as the text-based Lynx, still support FTP.
Syntax
FTP URL syntax is described in , taking the form: ftp://[user[:password]@]host[:port]/[url-path] (the bracketed parts are optional).
For example, the URL ftp://public.ftp-servers.example.com/mydirectory/myfile.txt represents the file myfile.txt from the directory mydirectory on the server public.ftp-servers.example.com as an FTP resource. The URL ftp://user001:secretpassword@private.ftp-servers.example.com/mydirectory/myfile.txt adds a specification of the username and password that must be used to access this resource.
More details on specifying a username and password may be found in the browsers' documentation (e.g., Firefox and Internet Explorer). By default, most web browsers use passive (PASV) mode, which more easily traverses end-user firewalls.
Some variation has existed in how different browsers treat path resolution in cases where there is a non-root home directory for a user.
Download manager
Most common download managers can receive files hosted on FTP servers, while some of them also give the interface to retrieve the files hosted on FTP servers. DownloadStudio allows not only download a file from FTP server but also view the list of files on a FTP server.
Other
LibreOffice declared its FTP support deprecated from 7.4 release, this was later removed in 24.2 release.
Security
FTP was not designed to be a secure protocol, and has many security weaknesses. In May 1999, the authors of listed a vulnerability to the following problems:
Brute-force attack
FTP bounce attack
Packet capture
Port stealing (guessing the next open port and usurping a legitimate connection)
Spoofing attack
Username enumeration
DoS or DDoS
FTP does not encrypt its traffic; all transmissions are in clear text, and usernames, passwords, commands and data can be read by anyone able to perform packet capture (sniffing) on the network. This problem is common to many of the Internet Protocol specifications (such as SMTP, Telnet, POP and IMAP) that were designed prior to the creation of encryption mechanisms such as TLS or SSL.
Common solutions to this problem include:
Using the secure versions of the insecure protocols, e.g., FTPS instead of FTP and TelnetS instead of Telnet.
Using a different, more secure protocol that can handle the job, e.g. SSH File Transfer Protocol or Secure Copy Protocol.
Using a secure tunnel such as Secure Shell (SSH) or virtual private network (VPN).
FTP over SSH
FTP over SSH is the practice of tunneling a normal FTP session over a Secure Shell connection. Because FTP uses multiple TCP connections (unusual for a TCP/IP protocol that is still in use), it is particularly difficult to tunnel over SSH. With many SSH clients, attempting to set up a tunnel for the control channel (the initial client-to-server connection on port 21) will protect only that channel; when data is transferred, the FTP software at either end sets up new TCP connections (data channels) and thus have no confidentiality or integrity protection.
Otherwise, it is necessary for the SSH client software to have specific knowledge of the FTP protocol, to monitor and rewrite FTP control channel messages and autonomously open new packet forwardings for FTP data channels. Software packages that support this mode include:
Tectia ConnectSecure (Win/Linux/Unix) of SSH Communications Security's software suite
FTP over SSH should not be confused with SSH File Transfer Protocol (SFTP).
Derivatives
FTPS
Explicit FTPS is an extension to the FTP standard that allows clients to request FTP sessions to be encrypted. This is done by sending the "AUTH TLS" command. The server has the option of allowing or denying connections that do not request TLS. This protocol extension is defined in . Implicit FTPS is an outdated standard for FTP that required the use of a SSL or TLS connection. It was specified to use different ports than plain FTP.
SSH File Transfer Protocol
The SSH file transfer protocol (chronologically the second of the two protocols abbreviated SFTP) transfers files and has a similar command set for users, but uses the Secure Shell protocol (SSH) to transfer files. Unlike FTP, it encrypts both commands and data, preventing passwords and sensitive information from being transmitted openly over the network. It cannot interoperate with FTP software, though some FTP client software offers support for the SSH file transfer protocol as well.
Trivial File Transfer Protocol
Trivial File Transfer Protocol (TFTP) is a simple, lock-step FTP that allows a client to get a file from or put a file onto a remote host. One of its primary uses is in the early stages of booting from a local area network, because TFTP is very simple to implement. TFTP lacks security and most of the advanced features offered by more robust file transfer protocols such as File Transfer Protocol. TFTP was first standardized in 1981 and the current specification for the protocol can be found in .
Simple File Transfer Protocol
Simple File Transfer Protocol (the first protocol abbreviated SFTP), as defined by , was proposed as an (unsecured) file transfer protocol with a level of complexity intermediate between TFTP and FTP. It was never widely accepted on the Internet, and is now assigned Historic status by the IETF. It runs through port 115, and often receives the initialism of SFTP. It has a command set of 11 commands and support three types of data transmission: ASCII, binary and continuous. For systems with a word size that is a multiple of 8 bits, the implementation of binary and continuous is the same. The protocol also supports login with user ID and password, hierarchical folders and file management (including rename, delete, upload, download, download with overwrite, and download with append).
FTP commands
FTP reply codes
Below is a summary of FTP reply codes that may be returned by an FTP server. These codes have been standardized in by the IETF. The reply code is a three-digit value. The first digit is used to indicate one of three possible outcomes — success, failure, or to indicate an error or incomplete reply:
2yz – Success reply
4yz or 5yz – Failure reply
1yz or 3yz – Error or Incomplete reply
The second digit defines the kind of error:
x0z – Syntax. These replies refer to syntax errors.
x1z – Information. Replies to requests for information.
x2z – Connections. Replies referring to the control and data connections.
x3z – Authentication and accounting. Replies for the login process and accounting procedures.
x4z – Not defined.
x5z – File system. These replies relay status codes from the server file system.
The third digit of the reply code is used to provide additional detail for each of the categories defined by the second digit.
| Technology | Internet | null |
53292 | https://en.wikipedia.org/wiki/Coca | Coca | Coca is any of the four cultivated plants in the family Erythroxylaceae, native to western South America. Coca is known worldwide for its psychoactive alkaloid, cocaine.
Different early-Holocene peoples in different areas of South America independently transformed Erythroxylum gracilipes plants into quotidian stimulant and medicinal crops now collectively called Coca. Archaeobotanical evidence show that Coca crops have been grown for well over 8,000 years in South America. They have had and still have a significant role in spiritual, economic, social and political dimensions for numerous indigenous cultures in the Andes and the Western Amazon arising from the use of the leaves as drugs and mild, daily stimulant.
The plant is grown as a cash crop in the Argentine Northwest, Bolivia, Alto Rio Negro Territory in Brazil, Colombia, Venezuela, Ecuador, and Peru, even in areas where its cultivation is unlawful. There are some reports that the plant is being cultivated in the south of Mexico, by using seeds imported from South America, as an alternative to smuggling its recreational product cocaine.
It also plays a fundamental role in many traditional Amazonian and Andean cultures as well as the Sierra Nevada de Santa Marta in northern Colombia.
The cocaine alkaloid content of dry Erythroxylum coca var. coca leaves was measured ranging from 0.23% to 0.96%. Coca-Cola used coca leaf extract in its products from 1885 until about 1903, when it began using decocainized leaf extract. Extraction of cocaine from coca requires several solvents and a chemical process known as an acid–base extraction, which can fairly easily extract the alkaloids from the plant.
Description
The coca plant resembles a blackthorn bush, and grows to a height of . The branches are curved, and the leaves are thin, opaque, oval, and taper at the extremities. A marked characteristic of the leaf is an areolated portion bounded by two longitudinal curved lines, one line on each side of the midrib, and more conspicuous on the under face of the leaf.
The flowers are small, and disposed in clusters on short stalks; the corolla is composed of five yellowish-white petals, the anthers are heart-shaped, and the pistil consists of three carpels united to form a three-chambered ovary. The flowers mature into red berries.
The leaves are sometimes eaten by the larvae of the moth Eloria noyesi.
Species and evolution
There are two species of coca crops, each with two varieties:
Erythroxylum coca
Erythroxylum coca var. coca (Bolivian or Huánuco Coca) – well adapted to the eastern Andes of Peru and Bolivia, an area of humid, tropical, montane forest.
Erythroxylum coca var. ipadu (Amazonian Coca) – cultivated in the lowland Amazon Basin in Peru and Colombia.
Erythroxylum novogranatense
Erythroxylum novogranatense var. novogranatense (Colombian Coca) – a highland variety that is utilized in lowland areas. It is cultivated in drier regions found in Colombia. However, E. novogranatense is very adaptable to varying ecological conditions. The leaves have parallel lines on either side of the central vein. These plants are called "Hayo" or "Ayu" among certain groups in Venezuela and Colombia.
Erythroxylum novogranatense var. truxillense (Trujillo Coca) – grown primarily in the Cajamarca and Amazonas states in Peru, including for the Empresa Nacional de la Coca S.A. and export by Coca-Cola for beverage flavoring.
All four of the cultivated cocas were domesticated from Erythroxylum gracilipes in pre-Columbian times, with significant archaeological sites reaching from Colombia to northern Chile, including the Las Vegas Culture in Ecuador, the Huaca Prieta site in Peru, and the Nanchoc valley in Peru – where leaf fragments and lime "cal" additives have been dated to over 8,000 years before present.
An initial theory of the origin and evolution of the cocas by Plowman and Bohm suggested that Erythroxylum coca var. coca is ancestral, while Erythroxylum novogranatense var. truxillense is derived from it to be drought tolerant, and Erythroxylum novogranatense var. novogranatense was further derived from Erythroxylum novogranatense var. truxillense in a linear series. In addition, E. coca var. ipadu was separately derived from E. coca var. coca when plants were taken into the Amazon basin.
Genetic evidence (Johnson et al. in 2005, Emche et al. in 2011, and Islam 2011) does not support this linear evolution. None of the four coca varieties are found in the wild, despite prior speculation by Plowman that wild populations of E. coca var. coca occur in the Huánuco and San Martín provinces of Peru. Recent phylogenetic evidence shows the closest wild relatives of the coca crops are Erythroxylum gracilipes Peyr. and Erythroxylum cataractarum Spruce ex. Peyr, and dense sampling of these species along with the coca crops from throughout their geographic ranges supports independent origins of domestication of Erythroxylum novogranatense and Erythroxylum coca from ancestor Erythroxylum gracilipes. It is possible that Amazonian coca was produced by yet a third independent domestication event from Erythroxylum gracilipes.
Thus, different early-Holocene peoples in different areas of South America independently transformed Erythroxylum gracilipes plants into quotidian stimulant and medicinal crops now collectively called coca.
Herbicide resistant varieties
Also known as supercoca or la millionaria, Boliviana negra is a relatively new form of coca that is resistant to a herbicide called glyphosate. Glyphosate is a key ingredient in the multibillion-dollar aerial coca eradication campaign undertaken by the government of Colombia with U.S. financial and military backing known as Plan Colombia.
The herbicide resistance of this strain has at least two possible explanations: that a "peer-to-peer" network of coca farmers used selective breeding to enhance this trait through tireless effort, or the plant was genetically modified in a laboratory. In 1996, a patented glyphosate-resistant soybean was marketed by Monsanto Company, suggesting that it would be possible to genetically modify coca in an analogous manner. Spraying Boliviana negra with glyphosate would serve to strengthen its growth by eliminating the non-resistant weeds surrounding it. Joshua Davis, in the Wired article cited below, found no evidence of CP4 EPSPS, a protein produced by the glyphosate-resistant soybean, suggesting Bolivana negra was either created in a lab by a different technique or bred in the field.
Cultivation
Coca is traditionally cultivated in the lower altitudes of the eastern slopes of the Andes (the Yungas), or the highlands depending on the species grown. Coca production begins in the valleys and upper jungle regions of the Andean region, where the countries of Colombia, Peru, and Bolivia are host to more than 98 percent of the global land area planted with coca. In the early 19th century, coca was cultivated in what is today the Dominican Republic (see Mayorasgo de Koka). In 2014, coca plantations were discovered in Mexico, and in 2020 in Honduras, which could have major implications for the illegal cultivation of the plant.
The seeds are sown from December to January in small plots () sheltered from the sun, and the young plants when at in height are placed in final planting holes (), or if the ground is level, in furrows () in carefully weeded soil. The plants thrive best in hot, damp and humid locations, such as the clearings of forests; but the leaves most preferred are obtained in drier areas, on the hillsides. The leaves are gathered from plants varying in age from one and a half to upwards of forty years, but only the new fresh growth is harvested. They are considered ready for plucking when they break on being bent. The first and most abundant harvest is in March after the rainy season, the second is at the end of June, and the third in October or November. The green leaves (matu) are spread in thin layers on coarse woollen cloths and dried in the sun; they are then packed in sacks, which must be kept dry in order to preserve the quality of the leaves.
Pharmacological aspects
The pharmacologically active ingredient of coca is the cocaine alkaloid, which has a concentration of about 0.3 to 1.5%, averaging 0.8%, in fresh leaves. Besides cocaine, the coca leaf contains a number of other alkaloids, including methylecgonine cinnamate, benzoylecgonine, truxilline, hydroxytropacocaine, tropacocaine, ecgonine, cuscohygrine, dihydrocuscohygrine, and hygrine. When chewed, coca acts as a mild stimulant and suppresses hunger, thirst, pain, and fatigue. Absorption of coca from the leaf is less rapid than nasal application of purified forms of the alkaloid (almost all of the coca alkaloid is absorbed within 20 minutes of nasal application, while it takes 2–12 hours after ingestion of the raw leaf for alkaline concentrations to peak.). When the raw leaf is consumed in tea, between 59 and 90% of the coca alkaloid is absorbed.
Coca users ingest between 60 and 80 milligrams of cocaine each time they chew the leaves according to United Nations Office on Drugs and Crime (UNODC). The coca leaf, when consumed in its natural form, does not induce a physiological or psychological dependence, nor does abstinence after long-term use produce symptoms typical to substance addiction. Due to its alkaloid content and non-addictive properties, coca has been suggested as a method to help recovering cocaine addicts to wean off the drug.
History
Traces of coca leaves found in northern Peru dates the communal chewing of coca with lime (the alkaline mineral, not the citrus fruit) 8,000 years back. Other evidence of coca traces have been found in mummies dating 3,000 years back in northern Chile. Beginning with the Valdivian culture, , there is an unbroken record of coca leaf consumption by succeeding cultural groups on the coast of Ecuador until European arrival as shown in their ceramic sculpture and abundant caleros or lime pots. Lime containers found in the north coast of Peru date around 2000 BC as evidenced by the findings at Huaca Prieta and the Jetetepeque river valley. Extensive archaeological evidence for the chewing of coca leaves dates back at least to the 6th century AD Moche period, and the subsequent Inca period, based on mummies found with a supply of coca leaves, pottery depicting the characteristic cheek bulge of a coca chewer, spatulas for extracting alkali and figured bags for coca leaves and lime made from precious metals, and gold representations of coca in special gardens of the Inca in Cuzco.
Coca chewing may originally have been limited to the eastern Andes before its introduction to the Inca. As the plant was viewed as having a divine origin, its cultivation became subject to a state monopoly and its use restricted to nobles and a few favored classes (court orators, couriers, favored public workers, and the army) by the rule of the Topa Inca (1471–1493). As the Incan empire declined, the leaf became more widely available. After some deliberation, Philip II of Spain issued a decree recognizing the drug as essential to the well-being of the Andean Indians but urging missionaries to end its religious use. The Spanish are believed to have effectively encouraged use of coca by an increasing majority of the population to increase their labor output and tolerance for starvation, but it is not clear that this was planned deliberately.
Andean people first started chewing coca leaf (Ertyhroxylum) and its popularity has been spread throughout the Northern and Central Andes, making its way down to Southern Central America, including areas like Bolivia, Chile, and Argentina. The coca leaf itself includes the active cocaine alkaloid which may be released through chewing or consumed in a powder-like form. This powder is usually extracted and made from burnt plant ashes, limestone or granite, and seashells. Andean people living in Central America have used a method to withdraw the lime from the coca plant using containers with sticks and have been able to indicate whether the coca leaves were either chewed historically even though many coca leaves haven't been discovered by archaeologists. There have been numerous effects that have been noted from the coca leaf as they are milder and more concentrated compared to pure cocaine. When Andean people began to first use the coca leaf, they noticed that it could produce a "high" and can be very addictive compared to tobacco if consumed in large quantities. Many Andean and Inca civilizations used to chew the coca leaf instead of consuming it as it provided a better "high" experience. Because of its strong addiction and high, the Incas only allowed this substance within honorary celebrations and rituals. Workers dealing with rigorous tasks such as long-distance travels, and more were allowed to take the substance as it eased their hardships along the way. There is little history before Andean people and the Incas to indicate if coca was restricted before these times and what instances it was initially used in. Sometimes coca leaves from the plant were used as offerings in rituals. Due to the nature of politics and religion in the Inca Empire, wealthy inhabitants handed out coca leaves during ritual ceremonies.
Coca was first introduced to Europe in the 16th century, but did not become popular until the mid-19th century, with the publication of an influential paper by Dr. Paolo Mantegazza praising its stimulating effects on cognition. This led to the invention of coca wine and the first production of pure cocaine. Coca wine (of which Vin Mariani was the best-known brand) and other coca-containing preparations were widely sold as patent medicines and tonics, with claims of a wide variety of health benefits. The original version of Coca-Cola was among these. These products became illegal in most countries outside of South America in the early 20th century, after the addictive nature of cocaine was widely recognized. In 1859, Albert Niemann of the University of Göttingen became the first person to isolate the chief alkaloid of coca, which he named "cocaine".
In the early 20th century, the Dutch colony of Java became a leading exporter of coca leaf. By 1912 shipments to Amsterdam, where the leaves were processed into cocaine, reached 1000 tons, overtaking the Peruvian export market. Apart from the years of the First World War, Java remained a greater exporter of coca than Peru until the end of the 1920s. Other colonial powers also tried to grow coca (including the British in India), but with the exception of the Japanese in Formosa, these were relatively unsuccessful.
In recent times (2006), the governments of several South American countries, such as Peru, Bolivia and Venezuela, have defended and championed the traditional use of coca, as well as the modern uses of the leaf and its extracts in household products such as teas and toothpaste. The coca plant was also the inspiration for Bolivia's Coca Museum.
Coca use by the Incas
Ethnohistorical sources
While many historians are in agreement that coca was a contributing factor to the daily life of the Inca, there are many different theories as to how this civilization came to adopt it as one of its staple crops and as a valued commodity. The Incas were able to accomplish significant things while stimulated by the effects of coca. The Incas did not have a graphical written language, but used the quipu, a fiber recording device. Spanish documents make it clear that coca was one of the most important elements of Inca culture. Coca was used in Inca feasts and religious rituals, among many other things. It was a driving factor in the labor efforts that Inca kings asked of their citizens, and also used to barter for other goods. Coca was vital to the Inca civilization and its culture. The Incas valued coca so much that they colonized tropical rain forests to the north and east of their capital in Cuzco so that they could increase and control their supply. The Incas colonized more humid regions because coca cannot grow above 2600 meters in elevation (coca is not frost-resistant).
Coca use in labor and military service
One of the most common uses of coca during the reign of the Inca was in the context of mit'a labor, a labor tax required of all able-bodied men in the Inca empire, and also in military service. Pedro Cieza de León wrote that the indigenous people of the Andes always seemed to have coca in their mouths. Mit'a laborers, soldiers, and others chewed coca to alleviate hunger and thirst while they were working and fighting. The results of this are evident in monumental construction and the successful expansion of the Inca empire through conquest. By chewing coca, laborers and soldiers were able to work harder and for longer periods. Some historians believe that coca and chicha (fermented corn beer) made it possible for the Incas to move large stones in order to create architectural masterpieces, especially ones of monolithic construction such as Sacsayhuamán.
Coca use in religious rituals
Due to the Spanish conquest of the Inca Empire, the Spaniards had direct access to the Inca. They had insight to their everyday lives, and it is through their lens that we learn about religion in the Inca Empire. While the indigenous author Pedro Cieza de León wrote about the effects coca had on the Inca, multiple Spanish men wrote about the importance of coca in their spirituality. For example Pedro Sarmiento de Gamboa, Father Bernabé Cobo, and Juan de Ulloa Mogollón noted how the Incas would leave coca leaves at important locations throughout the empire. They considered coca to be the highest form of plant offering that the Incas made.
The Incas would put coca leaves in the mouths of mummies, which were a sacred part of Inca culture. Mummies of Inca emperors were regarded for their wisdom and often consulted for important matters long after the body had deteriorated. Not only did many Inca mummies have coca leaves in their mouths, but they also carried coca leaves in bags. These are believed to be Inca sacrifices, and like the Aztecs, the Inca participated in sacrifices as well. It is clear that the Incas had a strong belief in the divinity of the coca leaf as there is now evidence that both the living and the dead were subjected to coca use. They even sent their sacrifices off to their death with a sacrificial bag of coca leaves. The coca leaf affected all stages of life for the Inca. Coca was also used in divination as ritual priests would burn a mixture of coca and llama fat and predict the future based on the appearance of the flame.
Coca use after the Spanish invasion and colonization
After the Spanish invasion and colonization of the Inca Empire, the use of coca was restricted and appropriated by the Spaniards. By many historical accounts, the Spaniards tried to eradicate the coca leaf from Inca life. The Spaniards enslaved Inca people and tried to prevent them from having "the luxury" of the coca leaf. Although the Spaniards noticed the state-controlled storage facilities that the Inca had built to distribute to its workers, they were still ignorant to plant spirit, divinity of coca, and the Incan admittance of the former. "This is my blood, this is my body" remembrance now was overshadowed by gates of behavior meeting efforts of worker control and service within work to spread concepts within outreach to support divinity and rights of the divine to exist in the divine's works. Not only that, enslaved Inca people were not capable of enduring the arduous labour the Spaniards made them do without using coca. Even though Spaniards were trying to push Catholicism onto the Inca, which did not allow them to eat before the Eucharist (the Spaniards thought coca to be food), they allowed them to continue to use coca to endure the labor associated with slavery. After seeing the effects and powers of the coca plant, many Spaniards saw another opportunity and started growing and selling coca themselves.
Traditional uses
Medicine
Traditional medical uses of coca are foremost as a stimulant to overcome fatigue, hunger, and thirst. It is considered particularly effective against altitude sickness. It also is used as an anesthetic and analgesic to alleviate the pain of headache, rheumatism, wounds and sores, etc. Before stronger anaesthetics were available, it also was used for broken bones, childbirth, and during trepanning operations on the skull. The high calcium content in coca explains why people used it for bone fractures. Because coca constricts blood vessels, it also serves to oppose bleeding, and coca seeds were used for nosebleeds. Indigenous use of coca has also been reported as a treatment for malaria, ulcers, asthma, to improve digestion, to guard against bowel laxity, as an aphrodisiac, and credited with improving longevity. Modern studies have supported a number of these medical applications.
Nutrition
Raw coca leaves, chewed or consumed as tea or mate de coca, are rich in nutritional properties. Specifically, the coca plant contains essential minerals (calcium, potassium, phosphorus), vitamins (B1, B2, C, and E) and nutrients such as protein and fiber.
Religion
Coca has also been a vital part of the religious cosmology of the Andean peoples of Peru, Chile, Bolivia, Ecuador, Colombia and northwest Argentina from the pre-Inca period through to the present. Coca leaves play a crucial part in offerings to the apus (mountains), Inti (the sun), or Pachamama (the earth). Coca leaves are also often read in a form of divination analogous to reading tea leaves in other cultures. As one example of the many traditional beliefs about coca, it is believed by the miners of Cerro de Pasco to soften the veins of ore, if masticated (chewed) and thrown upon them (see Cocamama in Inca mythology). In addition, coca use in shamanic rituals is well documented wherever local native populations have cultivated the plant. For example, the Tayronas of Colombia's Sierra Nevada de Santa Marta would chew the plant before engaging in extended meditation and prayer.
Chewing
In Bolivia bags of coca leaves are sold in local markets and by street vendors. The activity of chewing coca is called mambear, chacchar or acullicar, borrowed from Quechua, coquear (Northwest Argentina), or in Bolivia, picchar, derived from the Aymara language. The Spanish masticar is also frequently used, along with the slang term "bolear," derived from the word "bola" or ball of coca pouched in the cheek while chewing. Typical coca consumption varies between 20 and 60 grams per day, and contemporary methods are believed to be unchanged from ancient times. Coca is kept in a woven pouch (chuspa or huallqui). A few leaves are chosen to form a quid (acullico) held between the mouth and gums. Doing so may cause a tingling and numbing sensation in the mouth, in similar fashion to the formerly ubiquitous dental anaesthetic novocaine (as both cocaine and novocaine belong to the amino ester class of local anesthetics).
Chewing coca leaves is most common in indigenous communities across the central Andean region, particularly in places like the highlands of Argentina, Colombia, Bolivia, and Peru, where the cultivation and consumption of coca is a part of the national culture, similar to chicha. It also serves as a powerful symbol of indigenous cultural and religious identity, amongst a diversity of indigenous nations throughout South America. Chewing plants for medicinal mostly stimulating effects has a long history throughout the world: Khat in East Africa & the Arabian Peninsula, Tobacco in North America and Australia, and Areca nut in South/Southeast Asia & the Pacific Basin. Tobacco leaves were also traditionally chewed in the same way in North America (modern chewing tobacco is typically heavily processed). Khat chewing also has a history as a social custom dating back thousands of years analogous to the use of coca leaves.
One option for chewing coca is with a tiny quantity of ilucta (a preparation of the ashes of the quinoa plant) added to the coca leaves; it softens their astringent flavor and activates the alkaloids. Other names for this basifying substance are llipta in Peru and the Spanish word lejía, bleach in English. The consumer carefully uses a wooden stick (formerly often a spatula of precious metal) to transfer an alkaline component into the quid without touching his flesh with the corrosive substance. The alkali component, usually kept in a gourd (ishcupuro or poporo), can be made by burning limestone to form unslaked quicklime, burning quinoa stalks, or the bark from certain trees, and may be called llipta, tocra or mambe depending on its composition. Many of these materials are salty in flavor, but there are variations. The most common base in the La Paz area of Bolivia is a product known as lejía dulce (sweet lye), which is made from quinoa ashes mixed with aniseed and cane sugar, forming a soft black putty with a sweet and pleasing flavor. In some places, baking soda is used under the name bico.
In the Sierra Nevada de Santa Marta, on the Caribbean Coast of Colombia, coca is consumed by the Kogi, Arhuaco, and Wiwa by using a special device called poporo. The poporo is the mark of manhood; it is regarded by men as a good companion that means "food", "woman", "memory", and "meditation". When a boy is ready to be married, his mother initiates him in the use of the coca. This act of initiation is carefully supervised by the Mamo, a traditional priest-teacher-leader.
Fresh samples of the dried leaves, uncurled, are a deep green colour on the upper surface, and a grey-green on the lower surface, and have a strong tea-like aroma. When chewed, they produce a pleasurable numbness in the mouth, and have a pleasant, pungent taste. They are traditionally chewed with lime or some other reagent such as bicarbonate of soda to increase the release of the active ingredients from the leaf. Older species have a camphoraceous smell and a brownish color, and lack the pungent taste. | Biology and health sciences | Malpighiales | Plants |
53306 | https://en.wikipedia.org/wiki/Natural%20satellite | Natural satellite | A natural satellite is, in the most common usage, an astronomical body that orbits a planet, dwarf planet, or small Solar System body (or sometimes another natural satellite). Natural satellites are colloquially referred to as moons, a derivation from the Moon of Earth.
In the Solar System, there are six planetary satellite systems containing 288 known natural satellites altogether. Seven objects commonly considered dwarf planets by astronomers are also known to have natural satellites: , Pluto, Haumea, , Makemake, , and Eris. As of January 2022, there are 447 other minor planets known to have natural satellites.
A planet usually has at least around 10,000 times the mass of any natural satellites that orbit it, with a correspondingly much larger diameter. The Earth–Moon system is a unique exception in the Solar System; at 3,474 kilometres (2,158 miles) across, the Moon is 0.273 times the diameter of Earth and about of its mass. The next largest ratios are the Neptune–Triton system at 0.055 (with a mass ratio of about 1 to 4790), the Saturn–Titan system at 0.044 (with the second mass ratio next to the Earth–Moon system, 1 to 4220), the Jupiter–Ganymede system at 0.038, and the Uranus–Titania system at 0.031. For the category of dwarf planets, Charon has the largest ratio, being 0.52 the diameter and 12.2% the mass of Pluto.
Terminology
The first known natural satellite was the Moon, but it was considered a "planet" until Copernicus' introduction of De revolutionibus orbium coelestium in 1543. Until the discovery of the Galilean satellites in 1610 there was no opportunity for referring to such objects as a class. Galileo chose to refer to his discoveries as Planetæ ("planets"), but later discoverers chose other terms to distinguish them from the objects they orbited.
The first to use the term satellite to describe orbiting bodies was the German astronomer Johannes Kepler in his pamphlet Narratio de Observatis a se quatuor Iouis satellitibus erronibus ("Narration About Four Satellites of Jupiter Observed") in 1610. He derived the term from the Latin word satelles, meaning "guard", "attendant", or "companion", because the satellites accompanied their primary planet in their journey through the heavens.
The term satellite thus became the normal one for referring to an object orbiting a planet, as it avoided the ambiguity of "moon". In 1957, however, the launching of the artificial object Sputnik created a need for new terminology. The terms man-made satellite and artificial moon were very quickly abandoned in favor of the simpler satellite. As a consequence, the term has become linked primarily with artificial objects flown in space.
Because of this shift in meaning, the term moon, which had continued to be used in a generic sense in works of popular science and fiction, has regained respectability and is now used interchangeably with natural satellite, even in scientific articles. When it is necessary to avoid both the ambiguity of confusion with Earth's natural satellite the Moon and the natural satellites of the other planets on the one hand, and artificial satellites on the other, the term natural satellite (using "natural" in a sense opposed to "artificial") is used. To further avoid ambiguity, the convention is to capitalize the word Moon when referring to Earth's natural satellite (a proper noun), but not when referring to other natural satellites (common nouns).
Many authors define "satellite" or "natural satellite" as orbiting some planet or minor planet, synonymous with "moon" – by such a definition all natural satellites are moons, but Earth and other planets are not satellites.
A few recent authors define "moon" as "a satellite of a planet or minor planet", and "planet" as "a satellite of a star" – such authors consider Earth as a "natural satellite of the Sun".
Definition of a moon
There is no established lower limit on what is considered a "moon". Every natural celestial body with an identified orbit around a planet of the Solar System, some as small as a kilometer across, has been considered a moon, though objects a tenth that size within Saturn's rings, which have not been directly observed, have been called moonlets. Small asteroid moons (natural satellites of asteroids), such as Dactyl, have also been called moonlets.
The upper limit is also vague. Two orbiting bodies are sometimes described as a double planet rather than a primary and satellite. Asteroids such as 90 Antiope are considered double asteroids, but they have not forced a clear definition of what constitutes a moon. Some authors consider the Pluto–Charon system to be a double (dwarf) planet. The most common dividing line on what is considered a moon rests upon whether the barycentre is below the surface of the larger body, though this is somewhat arbitrary because it depends on distance as well as relative mass.
Origin and orbital characteristics
The natural satellites orbiting relatively close to the planet on prograde, uninclined circular orbits (regular satellites) are generally thought to have been formed out of the same collapsing region of the protoplanetary disk that created its primary. In contrast, irregular satellites (generally orbiting on distant, inclined, eccentric and/or retrograde orbits) are thought to be captured asteroids possibly further fragmented by collisions. Most of the major natural satellites of the Solar System have regular orbits, while most of the small natural satellites have irregular orbits. The Moon and the Moons of Pluto are exceptions among large bodies in that they are thought to have originated from the collision of two large protoplanetary objects early in the Solar System's history (see the giant impact hypothesis). The material that would have been placed in orbit around the central body is predicted to have reaccreted to form one or more orbiting natural satellites. As opposed to planetary-sized bodies, asteroid moons are thought to commonly form by this process. Triton is another exception; although large and in a close, circular orbit, its motion is retrograde and it is thought to be a captured dwarf planet.
Temporary satellites
The capture of an asteroid from a heliocentric orbit is not always permanent. According to simulations, temporary satellites should be a common phenomenon. The only observed examples are , , .
was a temporary satellite of Earth for nine months in 2006 and 2007.
Tidal locking
Most regular moons (natural satellites following relatively close and prograde orbits with small orbital inclination and eccentricity) in the Solar System are tidally locked to their respective primaries, meaning that the same side of the natural satellite always faces its planet. This phenomenon comes about through a loss of energy due to tidal forces raised by the planet, slowing the rotation of the satellite until it is negligible. Exceptions are known; one such exception is Saturn's natural satellite Hyperion, which rotates chaotically because of the gravitational influence of Titan. Pluto's four, circumbinary small moons also rotate chaotically due to Charon's influence.
In contrast, the outer natural satellites of the giant planets (irregular satellites) are too far away to have become locked. For example, Jupiter's Himalia, Saturn's Phoebe, and Neptune's Nereid have rotation periods in the range of ten hours, whereas their orbital periods are hundreds of days.
Satellites of satellites
No "moons of moons" or subsatellites (natural satellites that orbit a natural satellite of a planet) are currently known. In most cases, the tidal effects of the planet would make such a system unstable.
However, calculations performed after the 2008 detection of a possible ring system around Saturn's moon Rhea indicate that satellites orbiting Rhea could have stable orbits. Furthermore, the suspected rings are thought to be narrow, a phenomenon normally associated with shepherd moons. However, targeted images taken by the Cassini spacecraft failed to detect rings around Rhea.
It has also been proposed that Saturn's moon Iapetus had a satellite in the past; this is one of several hypotheses that have been put forward to account for its equatorial ridge.
Light-curve analysis suggests that Saturn's irregular satellite Kiviuq is extremely prolate, and is likely a contact binary or even a binary moon.
Trojan satellites
Two natural satellites are known to have small companions at both their and Lagrangian points, sixty degrees ahead and behind the body in its orbit. These companions are called trojan moons, as their orbits are analogous to the trojan asteroids of Jupiter. The trojan moons are Telesto and Calypso, which are the leading and following companions, respectively, of the Saturnian moon Tethys; and Helene and Polydeuces, the leading and following companions of the Saturnian moon Dione.
Asteroid satellites
The discovery of 243 Ida's natural satellite Dactyl in the early 1990s confirmed that some asteroids have natural satellites; indeed, 87 Sylvia has two. Some, such as 90 Antiope, are double asteroids with two comparably sized components.
Shape
Neptune's moon Proteus is the largest irregularly shaped natural satellite; the shapes of Eris' moon Dysnomia and ' moon Vanth are unknown. All other known natural satellites that are at least the size of Uranus's Miranda have lapsed into rounded ellipsoids under hydrostatic equilibrium, i.e. are "round/rounded satellites" and are sometimes categorized as planetary-mass moons. (Dysnomia's density is known to be high enough that it is probably a solid ellipsoid as well.) The larger natural satellites, being tidally locked, tend toward ovoid (egg-like) shapes: squat at their poles and with longer equatorial axes in the direction of their primaries (their planets) than in the direction of their motion. Saturn's moon Mimas, for example, has a major axis 9% greater than its polar axis and 5% greater than its other equatorial axis. Methone, another of Saturn's moons, is only around 3 km in diameter and visibly egg-shaped. The effect is smaller on the largest natural satellites, where their gravity is greater relative to the effects of tidal distortion, especially those that orbit less massive planets or, as in the case of the Moon, at greater distances.
Geological activity
Of the twenty known natural satellites in the Solar System that are large enough to be gravitationally rounded, several remain geologically active today. Io is the most volcanically active body in the Solar System, while Europa, Enceladus, Titan and Triton display evidence of ongoing tectonic activity and cryovolcanism. In the first three cases, the geological activity is powered by the tidal heating resulting from having eccentric orbits close to their giant-planet primaries. (This mechanism would have also operated on Triton in the past before its orbit was circularized.) Many other natural satellites, such as Earth's Moon, Ganymede, Tethys, and Miranda, show evidence of past geological activity, resulting from energy sources such as the decay of their primordial radioisotopes, greater past orbital eccentricities (due in some cases to past orbital resonances), or the differentiation or freezing of their interiors. Enceladus and Triton both have active features resembling geysers, although in the case of Triton solar heating appears to provide the energy. Titan and Triton have significant atmospheres; Titan also has hydrocarbon lakes. All four of the Galilean moons have atmospheres, though they are extremely thin. Four of the largest natural satellites, Europa, Ganymede, Callisto, and Titan, are thought to have subsurface oceans of liquid water, while smaller Enceladus also supports a global subsurface ocean of liquid water.
Occurrence in the Solar System
Besides planets and dwarf planets objects within our Solar System known to have natural satellites are 76 in the asteroid belt (five with two each), four Jupiter trojans, 39 near-Earth objects (two with two satellites each), and 14 Mars-crossers. There are also 84 known natural satellites of trans-Neptunian objects. Some 150 additional small bodies have been observed within the rings of Saturn, but only a few were tracked long enough to establish orbits. Planets around other stars are likely to have satellites as well, and although numerous candidates have been detected to date, none have yet been confirmed.
Of the inner planets, Mercury and Venus have no natural satellites; Earth has one large natural satellite, known as the Moon; and Mars has two tiny natural satellites, Phobos and Deimos.
The giant planets have extensive systems of natural satellites, including half a dozen comparable in size to Earth's Moon: the four Galilean moons, Saturn's Titan, and Neptune's Triton. Saturn has an additional six mid-sized natural satellites massive enough to have achieved hydrostatic equilibrium, and Uranus has five. It has been suggested that some satellites may potentially harbour life.
Among the objects generally agreed by astronomers to be dwarf planets, Ceres and have no known natural satellites. Pluto has the relatively large natural satellite Charon and four smaller natural satellites; Styx, Nix, Kerberos, and Hydra. Haumea has two natural satellites; , , Makemake, , and have one each. The Pluto–Charon system is unusual in that the center of mass lies in open space between the two, a characteristic sometimes associated with a double-planet system.
The seven largest natural satellites in the Solar System (those bigger than 2,500 km across) are Jupiter's Galilean moons (Ganymede, Callisto, Io, and Europa), Saturn's moon Titan, Earth's moon, and Neptune's captured natural satellite Triton. Triton, the smallest of these, has more mass than all smaller natural satellites together. Similarly in the next size group of nine mid-sized natural satellites, between 1,000 km and 1,600 km across, Titania, Oberon, Rhea, Iapetus, Charon, Ariel, Umbriel, Dione, and Tethys, the smallest, Tethys, has more mass than all smaller natural satellites together. As well as the natural satellites of the various planets, there are also over 80 known natural satellites of the dwarf planets, minor planets and other small Solar System bodies. Some studies estimate that up to 15% of all trans-Neptunian objects could have satellites.
The following is a comparative table classifying the natural satellites in the Solar System by diameter. The column on the right includes some notable planets, dwarf planets, asteroids, and trans-Neptunian objects for comparison. The natural satellites of the planets are named after mythological figures. These are predominantly Greek, except for the Uranian natural satellites, which are named after Shakespearean characters. The twenty satellites massive enough to be round are in bold in the table below. Minor planets and satellites where there is disagreement in the literature on roundness are italicized in the table below.
| Physical sciences | Astronomy | null |
53307 | https://en.wikipedia.org/wiki/Lepidoptera | Lepidoptera | Lepidoptera ( ) or lepidopterans is an order of winged insects which includes butterflies and moths. About 180,000 species of the Lepidoptera have been described, representing 10% of the total described species of living organisms, making it the second largest insect order (behind Coleoptera) with 126 families and 46 superfamilies, and one of the most widespread and widely recognizable insect orders in the world.
Lepidopteran species are characterized by more than three derived features. The most apparent is the presence of scales that cover the bodies, large triangular wings, and a proboscis for siphoning nectars. The scales are modified, flattened "hairs", and give butterflies and moths their wide variety of colors and patterns. Almost all species have some form of membranous wings, except for a few that have reduced wings or are wingless. Mating and the laying of eggs is normally performed near or on host plants for the larvae. Like most other insects, butterflies and moths are holometabolous, meaning they undergo complete metamorphosis. The larvae are commonly called caterpillars, and are completely different from their adult moth or butterfly forms, having a cylindrical body with a well-developed head, mandible mouth parts, three pairs of thoracic legs and from none up to five pairs of prolegs. As they grow, these larvae change in appearance, going through a series of stages called instars. Once fully matured, the larva develops into a pupa. A few butterflies and many moth species spin a silk casing or cocoon for protection prior to pupating, while others do not, instead going underground. A butterfly pupa, called a chrysalis, has a hard skin, usually with no cocoon. Once the pupa has completed its metamorphosis, a sexually mature adult emerges.
Lepidopterans first appeared in fossil record in the Triassic-Jurassic boundary and have coevolved with flowering plants since the angiosperm boom in the Middle/Late Cretaceous. They show many variations of the basic body structure that have evolved to gain advantages in lifestyle and distribution. Recent estimates suggest the order may have more species than earlier thought, and is among the five most species-rich orders (each with over 100,000 species) along with Coleoptera (beetles), Diptera (flies), Hymenoptera (ants, bees, wasps and sawflies) and Hemiptera (cicadas, aphids and other true bugs). They have, over millions of years, evolved a wide range of wing patterns and coloration ranging from drab moths akin to the related order Trichoptera, to the brightly colored and complex-patterned butterflies. Accordingly, this is the most recognized and popular of insect orders with many people involved in the observation, study, collection, rearing of, and commerce in these insects. A person who collects or studies this order is referred to as a lepidopterist.
Butterflies and moths are mostly herbivorous (folivorous) as caterpillars and nectarivorous as adults. They play an important role in the natural ecosystem as pollinators and serve as primary consumers in the food chain; conversely, their larvae (caterpillars) are considered very problematic to vegetation in agriculture, as they consume large quantity of plant matter (mostly foliage) to sustain growth. In many species, the female may produce from 200 to 600 eggs, while in others, the number may approach 30,000 eggs in one day. The caterpillars hatching from these eggs can cause significant damage to crops within a very short period of time. Many moth and butterfly species are of economic interest by virtue of their role as pollinators, the silk in their cocoon, or for extermination as pest species.
Etymology
The term Lepidoptera was used in 1746 by Carl Linnaeus in his Fauna Svecica. The word is derived from Greek , gen. ("scale") and ("wing"). Sometimes, the term Rhopalocera is used for the clade of all butterfly species, derived from the Ancient Greek () and () meaning "club" and "horn", respectively, coming from the shape of the antennae of butterflies.
The origins of the common names "butterfly" and "moth" are varied and often obscure. The English word butterfly is from Old English , with many variations in spelling. Other than that, the origin is unknown, although it could be derived from the pale yellow color of many species' wings suggesting the color of butter. The species of Heterocera are commonly called moths. The origins of the English word moth are clearer, deriving from Old English (cf. Northumbrian dialect ) from Common Germanic (compare Old Norse , Dutch and German all meaning "moth"). Perhaps its origins are related to Old English meaning "maggot" or from the root of "midge", which until the 16th century was used mostly to indicate the larva, usually in reference to devouring clothes.
The etymological origins of the word "caterpillar", the larval form of butterflies and moths, are from the early 16th century, from Middle English , , probably an alteration of Old North French (from Latin , "cat" + , "hairy").
Distribution and diversity
The Lepidoptera are among the most successful groups of insects. They are found on all continents, except Antarctica, and inhabit all terrestrial habitats ranging from desert to rainforest, from lowland grasslands to mountain plateaus, but almost always associated with higher plants, especially angiosperms (flowering plants). Among the most northern dwelling species of butterflies and moths is the Arctic Apollo (Parnassius arcticus), which is found in the Arctic Circle in northeastern Yakutia, at an altitude of above sea level. In the Himalayas, various Apollo species such as Parnassius epaphus have been recorded to occur up to an altitude of above sea level.
Some lepidopteran species exhibit symbiotic, phoretic, or parasitic lifestyles, inhabiting the bodies of organisms rather than the environment. Coprophagous pyralid moth species, called sloth moths, such as Bradipodicola hahneli and Cryptoses choloepi, are unusual in that they are exclusively found inhabiting the fur of sloths, mammals found in Central and South America. Two species of Tinea moths have been recorded as feeding on horny tissue and have been bred from the horns of cattle. The larva of Zenodochium coccivorella is an internal parasite of the coccid Kermes species. Many species have been recorded as breeding in natural materials or refuse such as owl pellets, bat caves, honeycombs or diseased fruit.
As of 2007, there were roughly 174,250 lepidopteran species described, with butterflies and skippers estimated to comprise around 17,950, and moths making up the rest. The vast majority of Lepidoptera are to be found in the tropics, but substantial diversity exists on most continents. North America has over 700 species of butterflies and over 11,000 species of moths, while about 400 species of butterflies and 14,000 species of moths are reported from Australia. The diversity of Lepidoptera in each faunal region has been estimated by John Heppner in 1991 based partly on actual counts from the literature, partly on the card indices in the Natural History Museum (London) and the National Museum of Natural History (Washington), and partly on estimates:
External morphology
Lepidoptera are morphologically distinguished from other orders principally by the presence of scales on the external parts of the body and appendages, especially the wings. Butterflies and moths vary in size from microlepidoptera only a few millimeters long, to conspicuous animals with a wingspan greater than , such as the Queen Alexandra's birdwing and Atlas moth.
Lepidopterans undergo a four-stage life cycle: egg; larva or caterpillar; pupa or chrysalis; and imago (plural: imagines) / adult and show many variations of the basic body structure, which give these animals advantages for diverse lifestyles and environments.
Head
The head is where many sensing organs and the mouth parts are found. Like the adult, the larva also has a toughened, or sclerotized head capsule. Here, two compound eyes, and chaetosema, raised spots or clusters of sensory bristles unique to Lepidoptera, occur, though many taxa have lost one or both of these spots. The antennae have a wide variation in form among species and even between different sexes. The antennae of butterflies are usually filiform and shaped like clubs, those of the skippers are hooked, while those of moths have flagellar segments variously enlarged or branched. Some moths have enlarged antennae or ones that are tapered and hooked at the ends.
The maxillary galeae are modified and form an elongated proboscis. The proboscis consists of one to five segments, usually kept coiled up under the head by small muscles when it is not being used to suck up nectar from flowers or other liquids. Some basal moths still have mandibles, or separate moving jaws, like their ancestors, and these form the family Micropterigidae.
The larvae, called caterpillars, have a toughened head capsule. Caterpillars lack the proboscis and have separate chewing mouthparts. These mouthparts, called mandibles, are used to chew up the plant matter that the larvae eat. The lower jaw, or labium, is weak, but may carry a spinneret, an organ used to create silk. The head is made of large lateral lobes, each having an ellipse of up to six simple eyes.
Thorax
The thorax is made of three fused segments, the prothorax, mesothorax, and metathorax, each with a pair of legs. The first segment contains the first pair of legs. In some males of the butterfly family Nymphalidae, the forelegs are greatly reduced and are not used for walking or perching. The three pairs of legs are covered with scales. Lepidoptera also have olfactory organs on their feet, which aid the butterfly in "tasting" or "smelling" out its food. In the larval form there are 3 pairs of true legs, with up to 11 pairs of abdominal legs (usually eight) and hooklets, called apical crochets.
The two pairs of wings are found on the middle and third segments, or mesothorax and metathorax, respectively. In the more recent genera, the wings of the second segment are much more pronounced, although some more primitive forms have similarly sized wings of both segments. The wings are covered in scales arranged like shingles, which form an extraordinary variety of colors and patterns. The mesothorax has more powerful muscles to propel the moth or butterfly through the air, with the wing of this segment (forewing) having a stronger vein structure. The largest superfamily, the Noctuoidea, has their wings modified to act as tympanal or hearing organs.
The caterpillar has an elongated, soft body that may have hair-like or other projections, three pairs of true legs, with none to 11 pairs of abdominal legs (usually eight) and hooklets, called apical crochets. The thorax usually has a pair of legs on each segment. The thorax is also lined with many spiracles on both the mesothorax and metathorax, except for a few aquatic species, which instead have a form of gills.
Abdomen
The abdomen, which is less sclerotized than the thorax, consists of 10 segments with membranes in between, allowing for articulated movement. The sternum, on the first segment, is small in some families and is completely absent in others. The last two or three segments form the external parts of the species' sex organs. The genitalia of Lepidoptera are highly varied and are often the only means of differentiating between species. Male genitals include a valva, which is usually large, as it is used to grasp the female during mating. Female genitalia include three distinct sections.
The females of basal moths have only one sex organ, which is used for copulation and as an ovipositor, or egg-laying organ. About 98% of moth species have a separate organ for mating, and an external duct that carries the sperm from the male.
The abdomen of the caterpillar has four pairs of prolegs, normally located on the third to sixth segments of the abdomen, and a separate pair of prolegs by the anus, which have a pair of tiny hooks called crotchets. These aid in gripping and walking, especially in species that lack many prolegs (e. g. larvae of Geometridae). In some basal moths, these prolegs may be on every segment of the body, while prolegs may be completely absent in other groups, which are more adapted to boring and living in sand (e. g., Prodoxidae and Nepticulidae, respectively).
Scales
The wings, head, and parts of the thorax and abdomen of Lepidoptera are covered with minute scales, a feature from which the order derives its name. Most scales are lamellar, or blade-like, and attached with a pedicel, while other forms may be hair-like or specialized as secondary sexual characteristics.
The lumen or surface of the lamella has a complex structure. It gives color either by colored pigments it contains, or through structural coloration with mechanisms that include photonic crystals and diffraction gratings.
Scales function in insulation, thermoregulation, producing pheromones (in males only), and aiding gliding flight, but the most important is the large diversity of vivid or indistinct patterns they provide, which help the organism protect itself by camouflage or mimicry, and which act as signals to other animals including rivals and potential mates.
Internal morphology
Reproductive system
In the reproductive system of butterflies and moths, the male genitalia are complex and unclear. In females the three types of genitalia are based on the relating taxa: 'monotrysian', 'exoporian', and 'ditrysian'. In the monotrysian type is an opening on the fused segments of the sterna 9 and 10, which act as insemination and oviposition. In the exoporian type (in Hepialoidea and Mnesarchaeoidea) are two separate places for insemination and oviposition, both occurring on the same sterna as the monotrysian type, i.e. 9 and 10. The ditrysian groups have an internal duct that carries sperm, with separate openings for copulation and egg-laying. In most species, the genitalia are flanked by two soft lobes, although they may be specialized and sclerotized in some species for ovipositing in area such as crevices and inside plant tissue. Hormones and the glands that produce them run the development of butterflies and moths as they go through their life cycles, called the endocrine system. The first insect hormone prothoracicotropic hormone (PTTH) operates the species life cycle and diapause. This hormone is produced by corpora allata and corpora cardiaca, where it is also stored. Some glands are specialized to perform certain task such as producing silk or producing saliva in the palpi. While the corpora cardiaca produce PTTH, the corpora allata also produces juvenile hormones, and the prothorocic glands produce moulting hormones.
Digestive system
In the digestive system, the anterior region of the foregut has been modified to form a pharyngeal sucking pump as they need it for the food they eat, which are for the most part liquids. An esophagus follows and leads to the posterior of the pharynx and in some species forms a form of crop. The midgut is short and straight, with the hindgut being longer and coiled. Ancestors of lepidopteran species, stemming from Hymenoptera, had midgut ceca, although this is lost in current butterflies and moths. Instead, all the digestive enzymes, other than initial digestion, are immobilized at the surface of the midgut cells. In larvae, long-necked and stalked goblet cells are found in the anterior and posterior midgut regions, respectively. In insects, the goblet cells excrete positive potassium ions, which are absorbed from leaves ingested by the larvae. Most butterflies and moths display the usual digestive cycle, but species with different diets require adaptations to meet these new demands. Some, like the luna moth, exhibit no digestive system whatsoever; they survive as adults from stored energy consumed as larvae and live for no longer than 7–10 days.
Circulatory system
In the circulatory system, hemolymph, or insect blood, is used to circulate heat in a form of thermoregulation, where muscles contraction produces heat, which is transferred to the rest of the body when conditions are unfavorable. In lepidopteran species, hemolymph is circulated through the veins in the wings by some form of pulsating organ, either by the heart or by the intake of air into the trachea.
Respiratory system
Air is taken in through spiracles along the sides of the abdomen and thorax supplying the trachea with oxygen as it goes through the lepidopteran's respiratory system. Three different tracheaes supply and diffuse oxygen throughout the species' bodies. The dorsal tracheae supply oxygen to the dorsal musculature and vessels, while the ventral tracheae supply the ventral musculature and nerve cord, and the visceral tracheae supply the guts, fat bodies, and gonads.
Polymorphism
Polymorphism is the appearance of forms or "morphs", which differ in color and number of attributes within a single species. In Lepidoptera, polymorphism can be seen not only between individuals in a population, but also between the sexes as sexual dimorphism, between geographically separated populations in geographical polymorphism, and between generations flying at different seasons of the year (seasonal polymorphism or polyphenism). In some species, the polymorphism is limited to one sex, typically the female. This often includes the phenomenon of mimicry when mimetic morphs fly alongside nonmimetic morphs in a population of a particular species. Polymorphism occurs both at specific level with heritable variation in the overall morphological adaptations of individuals, as well as in certain specific morphological or physiological traits within a species.
Environmental polymorphism, in which traits are not inherited, is often termed as polyphenism, which in Lepidoptera is commonly seen in the form of seasonal morphs, especially in the butterfly families of Nymphalidae and Pieridae. An Old World pierid butterfly, the common grass yellow (Eurema hecabe) has a darker summer adult morph, triggered by a long day exceeding 13 hours in duration, while the shorter diurnal period of 12 hours or less induces a paler morph in the postmonsoon period. Polyphenism also occurs in caterpillars, an example being the peppered moth, Biston betularia.
Geographical isolation causes a divergence of a species into different morphs. A good example is the Indian white admiral Limenitis procris, which has five forms, each geographically separated from the other by large mountain ranges. An even more dramatic showcase of geographical polymorphism is the Apollo butterfly (Parnassius apollo). Because the Apollos live in small local populations, thus having no contact with each other, coupled with their strong stenotopic nature and weak migration ability, interbreeding between populations of one species practically does not occur; by this, they form over 600 different morphs, with the size of spots on the wings of which varies greatly.
Sexual dimorphism is the occurrence of differences between males and females in a species. In Lepidoptera, it is widespread and almost completely set by genetic determination. Sexual dimorphism is present in all families of the Papilionoidea and more prominent in the Lycaenidae, Pieridae, and certain taxa of the Nymphalidae. Apart from color variation, which may differ from slight to completely different color-pattern combinations, secondary sexual characteristics may also be present. Different genotypes maintained by natural selection may also be expressed at the same time. Polymorphic and/or mimetic females occur in the case of some taxa in the Papilionidae primarily to obtain a level of protection not available to the male of their species. The most distinct case of sexual dimorphism is that of adult females of many Psychidae species which have only vestigial wings, legs, and mouthparts as compared to the adult males that are strong fliers with well-developed wings and feathery antennae.
Reproduction and development
Species of Lepidoptera undergo holometabolism or "complete metamorphosis". Their life cycle normally consists of an egg, a larva, a pupa, and an imago or adult. The larvae are commonly called caterpillars, and the pupae of moths encapsulated in silk are called cocoons, while the uncovered pupae of butterflies are called chrysalides.
Lepidopterans in diapause
Unless the species reproduces year-round, a butterfly or moth may enter diapause, a state of dormancy that allows the insect to survive unfavorable environmental conditions.
Mating
Males usually start eclosion (emergence) earlier than females and peak in numbers before females. Both of the sexes are sexually mature by the time of eclosion. Butterflies and moths normally do not associate with each other, except for migrating species, staying relatively asocial. Mating begins with an adult (female or male) attracting a mate, normally using visual stimuli, especially in diurnal species like most butterflies. However, the females of most nocturnal species, including almost all moth species, use pheromones to attract males, sometimes from long distances. Some species engage in a form of acoustic courtship, or attract mates using sound or vibration such as the polka-dot wasp moth, Syntomeida epilais.
Adaptations include undergoing one seasonal generation, two or even more, called voltinism (Univoltism, bivoltism, and multivism, respectively). Most lepidopterans in temperate climates are univoltine, while in tropical climates most have two seasonal broods. Some others may take advantage of any opportunity they can get, and mate continuously throughout the year. These seasonal adaptations are controlled by hormones, and these delays in reproduction are called diapause. Many lepidopteran species, after mating and laying their eggs, die shortly afterwards, having only lived for a few days after eclosion. Others may still be active for several weeks and then overwinter and become sexually active again when the weather becomes more favorable, or diapause. The sperm of the male that mated most recently with the female is most likely to have fertilized the eggs, but the sperm from a prior mating may still prevail.
Life cycle
Eggs
Lepidoptera usually reproduce sexually and are oviparous (egg-laying), though some species exhibit live birth in a process called ovoviviparity. A variety of differences in egg-laying and the number of eggs laid occur. Some species simply drop their eggs in flight (these species normally have polyphagous larvae, meaning they eat a variety of plants e. g., hepialids and some nymphalids) while most lay their eggs near or on the host plant on which the larvae feed. The number of eggs laid may vary from only a few to several thousand. The females of both butterflies and moths select the host plant instinctively, and primarily, by chemical cues.
The eggs are derived from materials ingested as a larva and in some species, from the spermatophores received from males during mating. An egg can only be 1/1000 the mass of the female, yet she may lay up to her own mass in eggs. Females lay smaller eggs as they age. Larger females lay larger eggs. The egg is covered by a hard-ridged protective outer layer of shell, called the chorion. It is lined with a thin coating of wax, which prevents the egg from drying out. Each egg contains a number of micropyles, or tiny funnel-shaped openings at one end, the purpose of which is to allow sperm to enter and fertilize the egg. Butterfly and moth eggs vary greatly in size between species, but they are all either spherical or ovate.
The egg stage lasts a few weeks in most butterflies, but eggs laid prior to winter, especially in temperate regions, go through diapause, and hatching may be delayed until spring. Other butterflies may lay their eggs in the spring and have them hatch in the summer. These butterflies are usually temperate species (e. g. Nymphalis antiopa).
Larvae
The larvae or caterpillars are the first stage in the life cycle after hatching. Caterpillars are "characteristic polypod larvae with cylindrical bodies, short thoracic legs, and abdominal prolegs (pseudopods)". They have a sclerotized head capsule with an adfrontal suture formed by medial fusion of the sclerites, mandibles (mouthparts) for chewing, and a soft tubular, segmented body, that may have hair-like or other projections, three pairs of true legs, and additional prolegs (up to five pairs). The body consists of thirteen segments, of which three are thoracic and ten are abdominal. Most larvae are herbivores, but a few are carnivores (some eat ants or other caterpillars) and detritivores.
Different herbivorous species have adapted to feed on every part of the plant and are normally considered pests to their host plants; some species have been found to lay their eggs on the fruit and other species lay their eggs on clothing or fur (e. g., Tineola bisselliella, the common clothes moth). Some species are carnivorous, and others are even parasitic. Some lycaenid species such as Phengaris rebeli are social parasites of Myrmica ant nests. A species of Geometridae from Hawaii has carnivorous larvae that catch and eat flies. Some pyralid caterpillars are aquatic.
The larvae develop rapidly with several generations in a year; however, some species may take up to 3 years to develop, and exceptional examples like Gynaephora groenlandica take as long as seven years. The larval stage is where the feeding and growing stages occur, and the larvae periodically undergo hormone-induced ecdysis, developing further with each instar, until they undergo the final larval-pupal molt.
The larvae of both butterflies and moths exhibit mimicry to deter potential predators. Some caterpillars have the ability to inflate parts of their heads to appear snake-like. Many have false eye-spots to enhance this effect. Some caterpillars have special structures called osmeteria (family Papilionidae), which are exposed to produce smelly chemicals used in defense. Host plants often have toxic substances in them, and caterpillars are able to sequester these substances and retain them into the adult stage. This helps make them unpalatable to birds and other predators. Such unpalatability is advertised using bright red, orange, black, or white warning colors. The toxic chemicals in plants are often evolved specifically to prevent them from being eaten by insects. Insects, in turn, develop countermeasures or make use of these toxins for their own survival. This "arms race" has led to the coevolution of insects and their host plants.
Wing development
No form of wing is externally visible on the larva, but when larvae are dissected, developing wings can be seen as disks, which can be found on the second and third thoracic segments, in place of the spiracles that are apparent on abdominal segments. Wing disks develop in association with a trachea that runs along the base of the wing, and are surrounded by a thin peripodial membrane, which is linked to the outer epidermis of the larva by a tiny duct. Wing disks are very small until the last larval instar, when they increase dramatically in size, are invaded by branching tracheae from the wing base that precede the formation of the wing veins and begin to develop patterns associated with several landmarks of the wing.
Near pupation, the wings are forced outside the epidermis under pressure from the hemolymph, and although they are initially quite flexible and fragile, by the time the pupa breaks free of the larval cuticle, they have adhered tightly to the outer cuticle of the pupa (in obtect pupae). Within hours, the wings form a cuticle so hard and well-joined to the body that pupae can be picked up and handled without damage to the wings.
Pupa
After about five to seven instars, or molts, certain hormones, like PTTH, stimulate the production of ecdysone, which initiates insect molting. The larva starts to develop into the pupa: body parts specific to the larva, such as the abdominal prolegs, degenerate, while others such as the legs and wings undergo growth. After finding a suitable place, the animal sheds its last larval cuticle, revealing the pupal cuticle underneath.
Depending on the species, the pupa may be covered in a silk cocoon, attached to different types of substrates, buried in the ground, or may not be covered at all. Features of the imago are externally recognizable in the pupa. All the appendages on the adult head and thorax are found cased inside the cuticle (antennae, mouthparts, etc.), with the wings wrapped around, adjacent to the antennae. The pupae of some species have functional mandibles, while the pupal mandibles are not functional in others.
Although the pupal cuticle is highly sclerotized, some of the lower abdominal segments are not fused, and are able to move using small muscles found in between the membrane. Moving may help the pupa, for example, escape the sun, which would otherwise kill it. The pupa of the Mexican jumping bean moth (Cydia saltitans) does this. The larvae cut a trapdoor in the bean (species of Sebastiania) and use the bean as a shelter. With a sudden rise in temperature, the pupa inside twitches and jerks, pulling on the threads inside. Wiggling may also help to deter parasitoid wasps from laying eggs on the pupa. Other species of moths are able to make clicks to deter predators.
The length of time before the pupa ecloses (emerges) varies greatly. The monarch butterfly may stay in its chrysalis for two weeks, while other species may need to stay for more than 10 months in diapause. The adult emerges from the pupa either by using abdominal hooks or from projections located on the head. The mandibles found in the most primitive moth families are used to escape from their cocoon (e. g., Micropterigoidea).
Adult
Most lepidopteran species do not live long after eclosion, only needing a few days to find a mate and then lay their eggs. Others may remain active for a longer period (from one to several weeks) or go through diapause and overwintering as monarch butterflies do, or waiting out environmental stress. Some adult species of microlepidoptera go through a stage where no reproductive-related activity occurs, lasting through summer and winter, followed by mating and oviposition in the early spring.
While most butterflies and moths are terrestrial, many species of Acentropinae (Crambidae) are truly aquatic with all stages except the adult occurring in water. Many species from other families such as Erebidae, Nepticulidae, Cosmopterigidae, Tortricidae, Olethreutidae, Noctuidae, Cossidae, and Sphingidae are aquatic or semiaquatic.
Video gallery of butterfly life cycle (Pieris rapae, the common cabbage white)
Behavior
Flight
Flight is an important aspect of the lives of butterflies and moths, and is used for evading predators, searching for food, and finding mates in a timely manner, as most lepidopteran species do not live long after eclosion. It is the main form of locomotion in most species. In Lepidoptera, the forewings and hindwings are mechanically coupled and flap in synchrony. Flight is anteromotoric, or being driven primarily by action of the forewings. Although lepidopteran species reportedly can still fly when their hindwings are cut off, it reduces their linear flight and turning capabilities.
Lepidopteran species have to be warm, about , to fly. They depend on their body temperature being sufficiently high and since they cannot regulate it themselves, this is dependent on their environment. Butterflies living in cooler climates may use their wings to warm their bodies. They will bask in the sun, spreading out their wings so that they get maximum exposure to the sunlight. In hotter climates butterflies can easily overheat, so they are usually active only during the cooler parts of the day, early morning, late afternoon or early evening. During the heat of the day, they rest in the shade. Some larger thick-bodied moths (e.g. Sphingidae) can generate their own heat to a limited degree by vibrating their wings. The heat generated by the flight muscles warms the thorax while the temperature of the abdomen is unimportant for flight. To avoid overheating, some moths rely on hairy scales, internal air sacs, and other structures to separate the thorax and abdomen and keep the abdomen cooler.
Some species of butterflies can reach fast speeds, such as the southern dart, which can go as fast as . Sphingids are some of the fastest flying insects, some are capable of flying at over , having a wingspan of . In some species, sometimes a gliding component to their flight exists. Flight occurs either as hovering, or as forward or backward motion. In butterfly and moth species, such as hawk moths, hovering is important as they need to maintain a certain stability over flowers when feeding on the nectar.
Navigation
Navigation is important to Lepidoptera species, especially for those that migrate. Butterflies, which have more species that migrate, have been shown to navigate using time-compensated sun compasses. They can see polarized light, so can orient even in cloudy conditions. The polarized light in the region close to the ultraviolet spectrum is suggested to be particularly important. Most migratory butterflies are those that live in semiarid areas where breeding seasons are short. The life histories of their host plants also influence the strategies of the butterflies. Other theories include the use of landscapes. Lepidoptera may use coastal lines, mountains, and even roads to orient themselves. Above sea, the flight direction is much more accurate if the coast is still visible.
Many studies have also shown that moths navigate. One study showed that many moths may use the Earth's magnetic field to navigate, as a study of the heart and dart moth suggests. Another study, of the migratory behavior of the silver Y, showed, even at high altitudes, the species can correct its course with changing winds, and prefers flying with favourable winds, suggesting a great sense of direction. Aphrissa statira in Panama loses its navigational capacity when exposed to a magnetic field, suggesting it uses the Earth's magnetic field.
Moths exhibit a tendency to circle artificial lights repeatedly. This suggests they use a technique of celestial navigation called transverse orientation. By maintaining a constant angular relationship to a bright celestial light, such as the Moon, they can fly in a straight line. Celestial objects are so far away, even after traveling great distances, the change in angle between the moth and the light source is negligible; further, the moon will always be in the upper part of the visual field or on the horizon. When a moth encounters a much closer artificial light and uses it for navigation, the angle changes noticeably after only a short distance, in addition to being often below the horizon. The moth instinctively attempts to correct by turning toward the light, causing airborne moths to come plummeting downwards, and at close range, which results in a spiral flight path that gets closer and closer to the light source. Other explanations have been suggested, such as the idea that moths may be impaired with a visual distortion called a Mach band by Henry Hsiao in 1972. He stated that they fly towards the darkest part of the sky in pursuit of safety, thus are inclined to circle ambient objects in the Mach band region.
Migration
Lepidopteran migration is typically seasonal, as the insects moving to escape dry seasons or other disadvantageous conditions. Most lepidopterans that migrate are butterflies, and the distance travelled varies. Some butterflies that migrate include the mourning cloak, painted lady, American lady, red admiral, and the common buckeye. A notable species of moth that migrates long distances is the bogong moth. The most well-known migrations are those of the eastern population of the monarch butterfly from Mexico to northern United States and southern Canada, a distance of about . Other well-known migratory species include the painted lady and several of the danaine butterflies. Spectacular and large-scale migrations associated with the monsoons are seen in peninsular India. Migrations have been studied in more recent times using wing tags and stable hydrogen isotopes.
Moths also undertake migrations, an example being the uraniids. Urania fulgens undergoes population explosions and massive migrations that may be not surpassed by any other insect in the Neotropics. In Costa Rica and Panama, the first population movements may begin in July and early August and depending on the year, may be very massive, continuing unabated for as long as five months.
Communication
Pheromones are commonly involved in mating rituals among species, especially moths, but they are also an important aspect of other forms of communication. Usually, the pheromones are produced by either the male or the female and detected by members of the opposite sex with their antennae. In many species, a gland between the eighth and ninth segments under the abdomen in the female produces the pheromones. Communication can also occur through stridulation, or producing sounds by rubbing various parts of the body together.
Moths are known to engage in acoustic forms of communication, most often as courtship, attracting mates using sound or vibration. Like most other insects, moths pick up these sounds using tympanic membranes in their abdomens. An example is that of the polka-dot wasp moth (Syntomeida epilais), which produces sounds with a frequency above that normally detectable by humans (about 20 kHz). These sounds also function as tactile communication, or communication through touch, as they stridulate, or vibrate a substrate like leaves and stems.
Most moths lack bright colors, as many species use coloration as camouflage, but butterflies engage in visual communication. Female cabbage butterflies, for example, use ultraviolet light to communicate, with scales colored in this range on the dorsal wing surface. When they fly, each down stroke of the wing creates a brief flash of ultraviolet light which the males apparently recognize as the flight signature of a potential mate. These flashes from the wings may attract several males that engage in aerial courtship displays.
Ecology
Moths and butterflies are important in the natural ecosystem. They are integral participants in the food chain; having co-evolved with flowering plants and predators, lepidopteran species have formed a network of trophic relationships between autotrophs and heterotrophs, which are included in the stages of Lepidoptera larvae, pupae, and adults. Larvae and pupae are links in the diets of birds and parasitic entomophagous insects. The adults are included in food webs in a much broader range of consumers (including birds, small mammals, reptiles, etc.).
Defense and predation
Lepidopteran species are soft bodied, fragile, and almost defenseless, while the immature stages move slowly or are immobile, hence all stages are exposed to predation. Adult butterflies and moths are preyed upon by birds, bats, lizards, amphibians, dragonflies, and spiders. One spider species, Argiope argentata, eats butterflies and moths and exhibits a long bite when preying on them rather than wrapping them in silk first. This is theorized to serve as an immobilization tactic. Caterpillars and pupae fall prey not only to birds, but also to invertebrate predators and small mammals, as well as fungi and bacteria. Parasitoid and parasitic wasps and flies may lay eggs in the caterpillar, which eventually kill it as they hatch inside its body and eat its tissues. Insect-eating birds are probably the largest predators. Lepidoptera, especially the immature stages, are an ecologically important food to many insectivorous birds, such as the great tit in Europe.
An "evolutionary arms race" can be seen between predator and prey species. The Lepidoptera have developed a number of strategies for defense and protection, including evolution of morphological characters and changes in ecological lifestyles and behaviors. These include aposematism, mimicry, camouflage, and development of threat patterns and displays. Only a few birds, such as the nightjars, hunt nocturnal lepidopterans. Their main predators are bats. Again, an "evolutionary race" exists, which has led to numerous evolutionary adaptations of moths to escape from their main predators, such as the ability to hear ultrasonic sounds, or even to emit sounds in some cases. Lepidopteran eggs are also preyed upon. Some caterpillars, such as the zebra swallowtail butterfly larvae, are cannibalistic.
Some species of Lepidoptera are poisonous to predators, such as the monarch butterfly in the Americas, Atrophaneura species (roses, windmills, etc.) in Asia, as well as Papilio antimachus, and the birdwings, the largest butterflies in Africa and Asia, respectively. They obtain their toxicity by sequestering the chemicals from the plants they eat into their own tissues. Some Lepidoptera manufacture their own toxins. Predators that eat poisonous butterflies and moths may become sick and vomit violently, learning not to eat those species. A predator which has previously eaten a poisonous lepidopteran may avoid other species with similar markings in the future, thus saving many other species, as well. Toxic butterflies and larvae tend to develop bright colors and striking patterns as an indicator to predators about their toxicity. This phenomenon is known as aposematism. Some caterpillars, especially members of Papilionidae, contain an osmeterium, a Y-shaped protrusible gland found in the prothoracic segment of the larvae. When threatened, the caterpillar emits unpleasant smells from the organ to ward off the predators.
Camouflage is also an important defense strategy, which involves the use of coloration or shape to blend into the surrounding environment. Some lepidopteran species blend with their surroundings, making them difficult to spot by predators. Caterpillars can exhibit shades of green that match its host plant. Caterpillars have been demonstrated to be able to detect the color of their surroundings and substrate using organs on their feet. Some caterpillars look like inedible objects, such as twigs or leaves. For instance, the mourning cloak fades into the backdrop of trees when it folds its wings back. The larvae of some species, such as the common Mormon (Papilio polytes) and the western tiger swallowtail look like bird droppings. For example, adult Sesiidae species (also known as clearwing moths) have a general appearance sufficiently similar to a wasp or hornet to make it likely the moths gain a reduction in predation by Batesian mimicry. Eyespots are a type of automimicry used by some butterflies and moths. In butterflies, the spots are composed of concentric rings of scales in different colors. The proposed role of the eyespots is to deflect attention of predators. Their resemblance to eyes provokes the predator's instinct to attack these wing patterns.
Batesian and Müllerian mimicry complexes are commonly found in Lepidoptera. Genetic polymorphism and natural selection give rise to otherwise edible species (the mimic) gaining a survival advantage by resembling inedible species (the model). Such a mimicry complex is referred to as Batesian and is most commonly known in the example between the limenitidine viceroy butterfly in relation to the inedible danaine monarch. The viceroy is, in fact, more toxic than the monarch and this resemblance should be considered as a case of Müllerian mimicry. In Müllerian mimicry, inedible species, usually within a taxonomic order, find it advantageous to resemble each other so as to reduce the sampling rate by predators that need to learn about the insects' inedibility. Taxa from the toxic genus Heliconius form one of the most well-known Müllerian complexes. The adults of the various species now resemble each other so well, the species cannot be distinguished without close morphological observation and, in some cases, dissection or genetic analysis.
Moths are able to hear the range emitted by bats, which in effect causes flying moths to make evasive maneuvers because bats are a main predator of moths. Ultrasonic frequencies trigger a reflex action in the noctuid moth that cause it to drop a few inches in its flight to evade attack. Tiger moths in a defense emit clicks within the same range of the bats, which interfere with the bats and foil their attempts to echolocate it.
Pollination
Most species of Lepidoptera engage in some form of entomophily (more specifically psychophily and phalaenophily for butterflies and moths, respectively), or the pollination of flowers. Most adult butterflies and moths feed on the nectar inside flowers, using their probosces to reach the nectar hidden at the base of the petals. In the process, the adults brush against the flowers' stamens, on which the reproductive pollen is made and stored. The pollen is transferred on appendages on the adults, which fly to the next flower to feed and unwittingly deposit the pollen on the stigma of the next flower, where the pollen germinates and fertilizes the seeds.
Flowers pollinated by butterflies tend to be large and flamboyant, pink or lavender in color, frequently having a landing area, and usually scented, as butterflies are typically day-flying. Since butterflies do not digest pollen (except for heliconid species,) more nectar is offered than pollen. The flowers have simple nectar guides, with the nectaries usually hidden in narrow tubes or spurs, reached by the long "tongue" of the butterflies. Butterflies such as Thymelicus flavus have been observed to engage in flower constancy, which means they are more likely to transfer pollen to other conspecific plants. This can be beneficial for the plants being pollinated, as flower constancy prevents the loss of pollen during different flights and the pollinators from clogging stigmas with pollen of other flower species.
Among the more important moth pollinator groups are the hawk moths of the family Sphingidae. Their behavior is similar to hummingbirds, i.e., using rapid wing beats to hover in front of flowers. Most hawk moths are nocturnal or crepuscular, so moth-pollinated flowers (e.g., Silene latifolia ) tend to be white, night-opening, large, and showy with tubular corollae and a strong, sweet scent produced in the evening, night, or early morning. A lot of nectar is produced to fuel the high metabolic rates needed to power their flight. Other moths (e.g., noctuids, geometrids, pyralids) fly slowly and settle on the flower. They do not require as much nectar as the fast-flying hawk moths, and the flowers tend to be small (though they may be aggregated in heads).
Mutualism
Mutualism is a form of biological interaction wherein each individual involved benefits in some way. An example of a mutualistic relationship would be that shared by yucca moths (Tegeculidae) and their host, yucca flowers (Asparagaceae). Female yucca moths enter the host flowers, collect the pollen into a ball using specialized maxillary palps, then move to the apex of the pistil, where pollen is deposited on the stigma, and lay eggs into the base of the pistil where seeds will develop. The larvae develop in the fruit pod and feed on a portion of the seeds. Thus, both insect and plant benefit, forming a highly mutualistic relationship. Another form of mutualism occurs between some larvae of butterflies and certain species of ants (e.g. Lycaenidae). The larvae communicate with the ants using vibrations transmitted through a substrate, such as the wood of a tree or stems, as well as using chemical signals. The ants provide some degree of protection to these larvae and they in turn gather honeydew secretions.
Parasitism
Only 42 species of parasitoid lepidopterans are known (1 Pyralidae; 40 Epipyropidae). The larvae of the greater and lesser wax moths feed on the honeycomb inside bee nests and may become pests; they are also found in bumblebee and wasp nests, albeit to a lesser extent. In northern Europe, the wax moth is regarded as the most serious parasitoid of the bumblebee and is found only in bumblebee nests. In some areas in southern England, as much as 80% of nests can be destroyed. Other parasitic larvae are known to prey upon cicadas and leaf hoppers.
In reverse, moths and butterflies may be subject to parasitic wasps and flies, which may lay eggs on the caterpillars, which hatch and feed inside its body, resulting in death. Although, in a form of parasitism called idiobiont, the adult paralyzes the host, so as not to kill it but for it to live as long as possible, for the parasitic larvae to benefit the most. In another form of parasitism, koinobiont, the species live off their hosts while inside (endoparasitic). These parasites live inside the host caterpillar throughout its life cycle or may affect it later on as an adult. In other orders, koinobionts include flies, a majority of coleopteran, and many hymenopteran parasitoids. Some species may be subject to a variety of parasites, such as the spongy moth (Lymantaria dispar), which is attacked by a series of 13 species, in six different taxa throughout its life cycle.
In response to a parasitoid egg or larva in the caterpillar's body, the plasmatocytes, or simply the host's cells can form a multilayered capsule that eventually causes the endoparasite to asphyxiate. The process, called encapsulation, is one of the caterpillar's only means of defense against parasitoids.
Other biological interactions
A few species of Lepidoptera are secondary consumers, or predators. These species typically prey upon the eggs of other insects, aphids, scale insects, or ant larvae. Some caterpillars are cannibals, and others prey on caterpillars of other species (e.g. Hawaiian Eupithecia ). Those of the 15 species in Eupithecia that mirror inchworms, are the only known species of butterflies and moths that are ambush predators. Four species are known to eat snails. For example, the Hawaiian caterpillar (Hyposmocoma molluscivora) uses silk traps, in a manner similar to that of spiders, to capture certain species of snails (typically Tornatellides).
Larvae of some species of moths in the Tineidae, Gelechiidae, and Noctuidae, besides others, feed on detritus, or dead organic material, such as fallen leaves and fruit, fungi, and animal products, and turn it into humus. Well-known species include the cloth moths (Tineola bisselliella, Tinea pellionella, and Trichophaga tapetzella), which feed on detritus containing keratin, including hair, feathers, cobwebs, bird nests (particularly of domestic pigeons, Columba livia domestica) and fruits or vegetables. These species are important to ecosystems as they remove substances that would otherwise take a long time to decompose.
In 2015 it was reported that wasp bracovirus DNA was present in Lepidoptera such as monarch butterflies, silkworms and moths. These were described in some newspaper articles as examples of a naturally occurring genetically engineered insects.
Evolution and systematics
History of study
Linnaeus in Systema Naturae (1758) recognized three divisions of the Lepidoptera: Papilio, Sphinx and Phalaena, with seven subgroups in Phalaena. These persist today as 9 of the superfamilies of Lepidoptera. Other works on classification followed including those by Michael Denis & Ignaz Schiffermüller (1775), Johan Christian Fabricius (1775) and Pierre André Latreille (1796). Jacob Hübner described many genera, and the lepidopteran genera were catalogued by Ferdinand Ochsenheimer and Georg Friedrich Treitschke in a series of volumes on the lepidopteran fauna of Europe published between 1807 and 1835. Gottlieb August Wilhelm Herrich-Schäffer (several volumes, 1843–1856), and Edward Meyrick (1895) based their classifications primarily on wing venation. Sir George Francis Hampson worked on the microlepidoptera during this period and Philipp Christoph Zeller published The Natural History of the Tineinae also on microlepidoptera (1855).
Among the first entomologists to study fossil insects and their evolution was Samuel Hubbard Scudder (1837–1911), who worked on butterflies. He published a study of the Florissant deposits of Colorado, including the exceptionally preserved Prodryas persephone. Andreas V. Martynov (1879–1938) recognized the close relationship between Lepidoptera and Trichoptera in his studies on phylogeny.
Major contributions in the 20th century included the creation of the monotrysia and ditrysia (based on female genital structure) by Borner in 1925 and 1939. Willi Hennig (1913–1976) developed the cladistic methodology and applied it to insect phylogeny. Niels P. Kristensen, E. S. Nielsen and D. R. Davis studied the relationships among monotrysian families and Kristensen worked more generally on insect phylogeny and higher Lepidoptera too. While it is often found that DNA-based phylogenies differ from those based on morphology, this has not been the case for the Lepidoptera; DNA phylogenies correspond to a large extent to morphology-based phylogenies.
Many attempts have been made to group the superfamilies of the Lepidoptera into natural groups, most of which fail because one of the two groups is not monophyletic: Microlepidoptera and Macrolepidoptera, Heterocera and Rhopalocera, Jugatae and Frenatae, Monotrysia and Ditrysia.
A 2024 genetic study found that the genomes of butterflies and moths have remained largely unchanged over the past 250 million years.
Fossil record
The fossil record for Lepidoptera is lacking in comparison to other winged species and tends not to be as common as some other insects in habitats that are most conducive to fossilization, such as lakes and ponds; their juvenile stage has only the head capsule as a hard part that might be preserved. Also, the scales covering their wings are hydrophobic and prevents their body from sinking when they end up on the water's surface. Lepidopteran bodies tend to come apart after death, and decompose quickly, so fossil remains are often extremely fragmentary. Of the fossils known, only an estimated 7% have been described. The location and abundance of the most common moth species are indicative that mass migrations of moths occurred over the Palaeogene North Sea, which is why there is a serious lack of moth fossils. Yet there are fossils, some preserved in amber and some in very fine sediments. Leaf mines are also seen in fossil leaves, although the interpretation of them is tricky.
Putative fossil stem group representatives of Amphiesmenoptera (the clade comprising Trichoptera and Lepidoptera) are known from the Triassic. The earliest known lepidopteran fossils are fossilized scales from the Triassic-Jurassic boundary. They were found as rare palynological elements in the sediments of the Triassic-Jurassic boundary from the cored Schandelah-1 well, drilled near Braunschweig in northern Germany. This pushes back the fossil record and origin of glossatan lepidopterans by about 70 million years, supporting molecular estimates of a Norian (ca 212 million years) divergence of glossatan and non-glossatan lepidopterans. The findings were reported in 2018 in the journal Science Advances. The authors of the study proposed that lepidopterans evolved a proboscis as an adaptation to drink from droplets and thin films of water for maintaining their fluid balance in the hot and arid climate of the Triassic.
The earliest named lepidopteran taxon is Archaeolepis mane, a primitive moth-like species from the Early Jurassic, dated back to around , and known only from three wings found in the Charmouth Mudstone of Dorset, UK. The wings show scales with parallel grooves under a scanning electron microscope and a characteristic wing venation pattern shared with Trichoptera (caddisflies). Only two more sets of Jurassic lepidopteran fossils have been found, as well as 13 sets from the Cretaceous, which all belong to primitive moth-like families.
Many more fossils are found from the Tertiary, and particularly the Eocene Baltic amber. The oldest genuine butterflies of the superfamily Papilionoidea have been found in the Paleocene MoClay or Fur Formation of Denmark. The best preserved fossil lepidopteran is the Eocene Prodryas persephone from the Florissant Fossil Beds.
Phylogeny
Lepidoptera and Trichoptera (caddisflies) are sister groups, sharing many similarities that are lacking in others; for example the females of both orders are heterogametic, meaning they have two different sex chromosomes, whereas in most species the males are heterogametic and the females have two identical sex chromosomes. The adults in both orders display a particular wing venation pattern on their forewings. The larvae in the two orders have mouth structures and glands with which they make and manipulate silk. Willi Hennig grouped the two orders into the superorder Amphiesmenoptera; together they are sister to the extinct order Tarachoptera. Lepidoptera descend from a diurnal moth-like common ancestor that either fed on dead or living plants.
The cladogram, based on molecular analysis, shows the order as a clade, sister to the Trichoptera, and more distantly related to the Diptera (true flies) and Mecoptera (scorpionflies).
The internal phylogeny of Lepidoptera is still being resolved. While many large clades have been established, interfamilial and superfamilial relationships are poorly understood. A large scale study by Regier et al. attempts to resolve these relationships using three different analysis methods, which is shown in the following cladogram.
Micropterigidae, Agathiphagidae and Heterobathmiidae are the oldest and most basal lineages of Lepidoptera. The adults of these families do not have the curled tongue or proboscis, that are found in most members of the order, but instead have chewing mandibles adapted for a special diet. Micropterigidae larvae feed on leaves, fungi, or liverworts (much like the Trichoptera). Adult Micropterigidae chew the pollen or spores of ferns. In the Agathiphagidae, larvae live inside kauri pines and feed on seeds. In Heterobathmiidae the larvae feed on the leaves of Nothofagus, the southern beech tree. These families also have mandibles in the pupal stage, which help the pupa emerge from the seed or cocoon after metamorphosis.
The Eriocraniidae have a short coiled proboscis in the adult stage, and though they retain their pupal mandibles with which they escaped the cocoon, their mandibles are non-functional thereafter. Most of these non-ditrysian families, are primarily leaf miners in the larval stage. In addition to the proboscis, there is a change in the scales among these basal lineages, with later lineages showing more complex perforated scales.
With the evolution of the Ditrysia in the mid-Cretaceous, there was a major reproductive change. The Ditrysia, which comprise 98% of the Lepidoptera, have two separate openings for reproduction in the females (as well as a third opening for excretion), one for mating, and one for laying eggs. The two are linked internally by a seminal duct. (In more basal lineages there is one cloaca, or later, two openings and an external sperm canal.) Of the early lineages of Ditrysia, Gracillarioidea and Gelechioidea are mostly leaf miners, but more recent lineages feed externally. In the Tineoidea, most species feed on plant and animal detritus and fungi, and build shelters in the larval stage.
The Yponomeutoidea is the first group to have significant numbers of species whose larvae feed on herbaceous plants, as opposed to woody plants. They evolved about the time that flowering plants underwent an expansive adaptive radiation in the mid-Cretaceous, and the Gelechioidea that evolved at this time also have great diversity. Whether the processes involved coevolution or sequential evolution, the diversity of the Lepidoptera and the angiosperms increased together.
In the so-called "macrolepidoptera", which constitutes about 60% of lepidopteran species, there was a general increase in size, better flying ability (via changes in wing shape and linkage of the forewings and hindwings), reduction in the adult mandibles, and a change in the arrangement of the crochets (hooks) on the larval prolegs, perhaps to improve the grip on the host plant. Many also have tympanal organs, that allow them to hear. These organs evolved eight times, at least, because they occur on different body parts and have structural differences.
The main lineages in the macrolepidoptera are the Noctuoidea, Bombycoidea, Lasiocampidae, Mimallonoidea, Geometroidea and Rhopalocera. Bombycoidea plus Lasiocampidae plus Mimallonoidea may be a monophyletic group. The Rhopalocera, comprising the Papilionoidea (butterflies), Hesperioidea (skippers), and the Hedyloidea (moth-butterflies), are the most recently evolved. There is quite a good fossil record for this group, with the oldest skipper dating from .
Taxonomy
Taxonomy is the classification of species in selected taxa, the process of naming being called nomenclature. There are over 120 families in Lepidoptera, in 45 to 48 superfamilies. Lepidoptera have always been, historically, classified in five suborders, one of which is of primitive moths that never lost the morphological features of their ancestors. The rest of the moths and butterflies make up ninety-eight percent of the other taxa, making Ditrysia. More recently, findings of new taxa, larvae and pupa have aided in detailing the relationships of primitive taxa, phylogenetic analysis showing the primitive lineages to be paraphyletic compared to the rest of Lepidoptera lineages. Recently, lepidopterists have abandoned clades like suborders, and those between orders and superfamilies.
Zeugloptera is a clade with Micropterigoidea being its only superfamily, containing the single family Micropterigidae. Species of Micropterigoidea are practically living fossils, being one of the most primitive lepidopteran groups, still retaining chewing mouthparts (mandibles) in adults, unlike other clades of butterflies and moths. About 120 species are known worldwide, with more than half the species in the genus Micropterix in the Palearctic region. There are only two known in North America (Epimartyria), with many more being found in Asia and the southwest Pacific, particularly New Zealand with about 50 species.
Aglossata is the second most primitive lineage of Lepidoptera; being first described in 1952 by Lionel Jack Dumbleton. Agathiphagidae is the only family in Aglossata and contains two species in its only genus, Agathiphaga. Agathiphaga queenslandensis and Agathiphaga vitiensis are found along the north-eastern coast of Queensland, Australia, and in Fiji to Vanuatu and the Solomon Islands, respectively.
Heterobathmiina was first described by Kristensen and Nielsen in 1979. Heterobathmiidae is the only family and includes about 10 species, which are day-flying, metallic moths, confined to southern South America, the adults eat the pollen of Nothofagus or southern beech and the larvae mine the leaves.
Glossata contains a majority of the species, with the most obvious difference being non-functioning mandibles, and elongated maxillary galeae or the proboscis. The basal clades still retaining some of the ancestral features of the wings such as similarly shaped fore- and hindwings with relatively complete venation. Glossata also contains the division Ditrysia, which contains 98% of all described species in Lepidoptera.
Relationship to people
Culture
Artistic depictions of butterflies have been used in many cultures including as early as 3500 years ago, in Egyptian hieroglyphs. Today, butterflies are widely used in various objects of art and jewelry: mounted in frames, embedded in resin, displayed in bottles, laminated in paper, and in some mixed media artworks and furnishings. Butterflies have also inspired the "butterfly fairy" as an art and fictional character.
In many cultures the soul of a dead person is associated with the butterfly, for example in Ancient Greece, where the word for butterfly ψυχή (psyche) also means soul and breath. In Latin, as in Ancient Greece, the word for "butterfly" papilio was associated with the soul of the dead. The skull-like marking on the thorax of the death's-head hawkmoth has helped these moths, particularly A. atropos, earn a negative reputation, such as associations with the supernatural and evil. The moth has been prominently featured in art and movies such as (by Buñuel and Dalí) and The Silence of the Lambs, and in the artwork of the Japanese metal band Sigh's album Hail Horror Hail. According to Kwaidan: Stories and Studies of Strange Things, by Lafcadio Hearn, a butterfly was seen in Japan as the personification of a person's soul; whether they be living, dying, or already dead. One Japanese superstition says that if a butterfly enters your guestroom and perches behind the bamboo screen, the person whom you most love is coming to see you. However, large numbers of butterflies are viewed as bad omens. When Taira no Masakado was secretly preparing for his famous revolt, there appeared in Kyoto so vast a swarm of butterflies that the people were frightened—thinking the apparition to be a portent of coming evil.
In the ancient Mesoamerican city of Teotihuacan, the brilliantly colored image of the butterfly was carved into many temples, buildings, jewelry, and emblazoned on incense burners in particular. The butterfly was sometimes depicted with the maw of a jaguar and some species were considered to be the reincarnations of the souls of dead warriors. The close association of butterflies to fire and warfare persisted through to the Aztec civilization and evidence of similar jaguar-butterfly images has been found among the Zapotec, and Maya civilizations.
Pests
The larvae of many lepidopteran species are major pests in agriculture. Some of the major pests include Tortricidae, Noctuidae, and Pyralidae. The larvae of the Noctuidae genus Spodoptera (armyworms), Helicoverpa (corn earworm), or Pieris brassicae can cause extensive damage to certain crops. Helicoverpa zea larvae (cotton bollworms or tomato fruitworms) are polyphagous, meaning they eat a variety of crops, including tomatoes and cotton. Peridroma saucia (variegated cutworms) are described as one of the most damaging pests to gardens, with the ability to destroy entire gardens and fields in a matter of days.
Butterflies and moths are one of the largest taxa to solely feed and be dependent on living plants, in terms of the number of species, and they are in many ecosystems, making up the largest biomass to do so. In many species, the female may produce anywhere from 200 to 600 eggs, while in some others it may go as high as 30,000 eggs in one day. This can create many problems for agriculture, where many caterpillars can affect acres of vegetation. Some reports estimate that there have been over 80,000 caterpillars of several different taxa feeding on a single oak tree. In some cases, phytophagous larvae can lead to the destruction of entire trees in relatively short periods of time.
Ecological ways of removing pest Lepidoptera species are becoming more economically viable, as research has shown ways like introducing parasitic wasps and flies. For example, Sarcophaga aldrichi, a fly which deposited larvae feed upon the pupae of the forest tent caterpillar moth. Pesticides can affect other species other than the species they are targeted to eliminate, damaging the natural ecosystem. Another good biological pest control method is the use of pheromone traps. A pheromone trap is a type of insect trap that uses pheromones to lure insects. Sex pheromones and aggregating pheromones are the most common types used. A pheromone-impregnated lure is encased in a conventional trap such as a Delta trap, water-pan trap, or funnel trap.
Species of moths that are detritivores would naturally eat detritus containing keratin, such as hairs or feathers. Well known species are cloth moths (T. bisselliella, T. pellionella, and T. tapetzella), feeding on foodstuffs that people find economically important, such as cotton, linen, silk and wool fabrics as well as furs; furthermore they have been found on shed feathers and hair, bran, semolina and flour (possibly preferring wheat flour), biscuits, casein, and insect specimens in museums.
Beneficial insects
Even though some butterflies and moths affect the economy negatively, many species are a valuable economic resource. The most prominent example is that of the domesticated silkworm moth (Bombyx mori), the larvae of which make their cocoons out of silk, which can be spun into cloth. Silk is and has been an important economic resource throughout history. The species Bombyx mori has been domesticated to the point where it is completely dependent on mankind for survival. A number of wild moths such as Bombyx mandarina, and Antheraea species, besides others, provide commercially important silks.
The preference of the larvae of most lepidopteran species to feed on a single species or limited range of plants is used as a mechanism for biological control of weeds in place of herbicides. The pyralid cactus moth was introduced from Argentina to Australia, where it successfully suppressed millions of acres of prickly pear cactus. Another species of the Pyralidae, called the alligator weed stem borer (Arcola malloi), was used to control the aquatic plant known as alligator weed (Alternanthera philoxeroides) in conjunction with the alligator weed flea beetle; in this case, the two insects work in synergy and the weed rarely recovers.
Breeding butterflies and moths, or butterfly gardening/rearing, has become an ecologically viable process of introducing species into the ecosystem to benefit it. Butterfly ranching in Papua New Guinea permits nationals of that country to "farm" economically valuable insect species for the collectors market in an ecologically sustainable manner.
Food
Lepidoptera feature prominently in entomophagy as food items on almost every continent. While in most cases, adults, larvae or pupae are eaten as staples by indigenous people, beondegi or silkworm pupae are eaten as a snack in Korean cuisine while Maguey worm is considered a delicacy in Mexico. In some parts of Huasteca, the silk nests of the Madrone butterfly are maintained on the edge of roof tops of houses for consumption. In the Carnia region of Italy, children catch and eat ingluvies of the toxic Zygaena moths in early summer. The ingluvies, despite having a very low cyanogenic content, serve as a convenient, supplementary source of sugar to the children who can include this resource as a seasonal delicacy at minimum risk. Outside of this instance, adult Lepidoptera are rarely consumed by humans, with the sole exception of the Bogong moth.
Health
Some larvae of both moths and butterflies have a form of hair that has been known to be a cause of human health problems. Caterpillar hairs sometimes have toxins in them and species from approximately 12 families of moths or butterflies worldwide can inflict serious human injuries (urticarial dermatitis and atopic asthma to osteochondritis, consumption coagulopathy, renal failure, and intracerebral hemorrhage). Skin rashes are the most common, but there have been fatalities. Lonomia is a frequent cause of envenomation in humans in Brazil, with 354 cases reported between 1989 and 2005. Lethality ranging up to 20% with death caused most often by intracranial hemorrhage.
These hairs have also been known to cause keratoconjunctivitis. The sharp barbs on the end of caterpillar hairs can get lodged in soft tissues and mucous membranes such as the eyes. Once they enter such tissues, they can be difficult to extract, often exacerbating the problem as they migrate across the membrane. This becomes a particular problem in an indoor setting. The hairs easily enter buildings through ventilation systems and accumulate in indoor environments because of their small size, which makes it difficult for them to be vented out. This accumulation increases the risk of human contact in indoor environments.
| Biology and health sciences | Lepidoptera | null |
53309 | https://en.wikipedia.org/wiki/Psychedelic%20drug | Psychedelic drug | Psychedelics are a subclass of hallucinogenic drugs whose primary effect is to trigger non-ordinary mental states (known as psychedelic experiences or "trips") and a perceived "expansion of consciousness". Also referred to as classic hallucinogens or serotonergic hallucinogens, the term psychedelic is sometimes used more broadly to include various types of hallucinogens, such as those which are atypical or adjacent to psychedelia like salvia and MDMA, respectively.
Classic psychedelics generally cause specific psychological, visual, and auditory changes, and oftentimes a substantially altered state of consciousness. They have had the largest influence on science and culture, and include mescaline, LSD, psilocybin, and DMT.
Most psychedelic drugs fall into one of the three families of chemical compounds: tryptamines, phenethylamines, or lysergamides (LSD is considered both a tryptamine and lysergamide). They act via serotonin 2A receptor agonism. When compounds bind to serotonin 5-HT2A receptors, they modulate the activity of key circuits in the brain involved with sensory perception and cognition. However, the exact nature of how psychedelics induce changes in perception and cognition via the 5-HT2A receptor is still unknown. The psychedelic experience is often compared to non-ordinary forms of consciousness such as those experienced in meditation, mystical experiences, and near-death experiences, which also appear to be partially underpinned by altered default mode network activity. The phenomenon of ego death is often described as a key feature of the psychedelic experience.
Many psychedelic drugs are illegal to possess without lawful authorisation, exemption or license worldwide under the UN conventions, with occasional exceptions for religious use or research contexts. Despite these controls, recreational use of psychedelics is common. Legal barriers have made the scientific study of psychedelics more difficult. Research has been conducted, however, and studies show that psychedelics are physiologically safe and rarely lead to addiction. Studies conducted using psilocybin in a psychotherapeutic setting reveal that psychedelic drugs may assist with treating depression, alcohol addiction, and nicotine addiction. Although further research is needed, existing results suggest that psychedelics could be effective treatments for certain forms of psychopathology. A 2022 survey found that 28% of Americans had used a psychedelic at some point in their life.
Etymology and nomenclature
The term psychedelic was coined by the psychiatrist Humphrey Osmond during written correspondence with author Aldous Huxley (written in a rhyme: “To fathom Hell or soar angelic/Just take a pinch of psychedelic.”) and presented to the New York Academy of Sciences by Osmond in 1957. It is irregularly derived from the Greek words ψυχή (psychḗ, meaning 'mind, soul') and δηλείν (dēleín, meaning 'to manifest'), with the intended meaning "mind manifesting" or alternatively "soul manifesting", and the implication that psychedelics can reveal unused potentials of the human mind. The term was loathed by American ethnobotanist Richard Schultes but championed by American psychologist Timothy Leary.
Aldous Huxley had suggested his own coinage phanerothyme (Greek phaneroein- "to make manifest or visible" and Greek thymos "soul", thus "to reveal the soul") to Osmond in 1956. Recently, the term entheogen (meaning "that which produces the divine within") has come into use to denote the use of psychedelic drugs, as well as various other types of psychoactive substances, in a religious, spiritual, and mystical context.
In 2004, David E. Nichols wrote the following about the nomenclature used for psychedelic drugs:
Robin Carhart-Harris and Guy Goodwin write that the term psychedelic is preferable to hallucinogen for describing classical psychedelics because of the term hallucinogens "arguably misleading emphasis on these compounds' hallucinogenic properties."
While the term psychedelic is most commonly used to refer only to serotonergic hallucinogens, it is sometimes used for a much broader range of drugs, including empathogen–entactogens, dissociatives, and atypical hallucinogens/psychoactives such as Amanita muscaria, Cannabis sativa, Nymphaea nouchali and Salvia divinorum. Thus, the term serotonergic psychedelic is sometimes used for the narrower class. It is important to check the definition of a given source. This article uses the more common, narrower definition of psychedelic.
Examples
2C-B (2,5-dimethoxy-4-bromophenethylamine) is a substituted phenethylamine first synthesised in 1974 by Alexander Shulgin. 2C-B is both a psychedelic and a mild entactogen, with its psychedelic effects increasing and its entactogenic effects decreasing with dosage. 2C-B is the most well known compound in the 2C family, their general structure being discovered as a result of modifying the structure of mescaline.
DMT (N,N-dimethyltryptamine) is an indole alkaloid found in various species of plants. Traditionally it is consumed by tribes in South America in the form of ayahuasca. A brew is used that consists of DMT-containing plants as well as plants containing MAOIs, specifically harmaline, which allows DMT to be consumed orally without being rendered inactive by monoamine oxidase enzymes in the digestive system. In the Western world DMT is more commonly consumed via the vaporisation of freebase DMT. Whereas Ayahuasca typically lasts for several hours, inhalation has an onset measured in seconds and has effects measured in minutes, being significantly more intense. Particularly in vaporised form, DMT has the ability to cause users to enter a hallucinatory realm fully detached from reality, being typically characterised by hyperbolic geometry, and described as defying visual or verbal description. Users have also reported encountering and communicating with entitites within this hallucinatory state. DMT is the archetypal substituted tryptamine, being the structural scaffold of psilocybin and – to a lesser extent – the lysergamides.
LSD (Lysergic acid diethylamide) is a derivative of lysergic acid, which is obtained from the hydrolysis of ergotamine. Ergotamine is an alkaloid found in the fungus Claviceps purpurea, which primarily infects rye. LSD is both the prototypical psychedelic and the prototypical lysergamide. As a lysergamide, LSD contains both a tryptamine and phenethylamine group within its structure. As a result of containing a phenethylamine group LSD agonises dopamine receptors as well as serotonin receptors, making it more energetic in effect in contrast to the more sedating effects of psilocin, which is not a dopamine agonist.
Mescaline (3,4,5-trimethoxyphenethylamine) is a phenethylamine alkaloid found in various species of cacti, the best-known of these being peyote (Lophophora williamsii) and the San Pedro cactus (Trichocereus macrogonus var. pachanoi, syn. Echinopsis pachanoi). Mescaline has effects comparable to those of LSD and psilocybin, albeit with a greater emphasis on colors and patterns. Ceremonial San Pedro use seems to be characterized by relatively strong spiritual experiences, and low incidence of challenging experiences.
Psilocin (4-HO-DMT) is the dephosphorylated active metabolite of the indole alkaloid psilocybin and a substituted tryptamine, which is produced in over 200 species of fungi. Of the Classical psychedelics psilocybin has attracted the greatest academic interest regarding its ability to manifest mystical experiences, although all psychedelics are capable of doing so to variable degrees. O-Acetylpsilocin (4-AcO-DMT) is an acetylated analog of psilocin. Additionally, replacement of a methyl group at the dimethylated nitrogen with an isopropyl or ethyl group yields 4-HO-MIPT and 4-HO-MET, respectively.
Uses
Traditional
A number of frequently mentioned or traditional psychedelics such as Ayahuasca (which contains DMT), San Pedro, Peyote, and Peruvian torch (which all contain mescaline), Psilocybe mushrooms (which contain psilocin/psilocybin) and Tabernanthe iboga (which contains the unique psychedelic ibogaine) all have a long and extensive history of spiritual, shamanic and traditional usage by indigenous peoples in various world regions, particularly in Latin America, but also Gabon, Africa in the case of iboga. Different countries and/or regions have come to be associated with traditional or spiritual use of particular psychedelics, such as the ancient and entheogenic use of psilocybe mushrooms by the native Mazatec people of Oaxaca, Mexico or the use of the ayahuasca brew in the Amazon basin, particularly in Peru for spiritual and physical healing as well as for religious festivals. Peyote has also been used for several thousand years in the Rio Grande Valley in North America by native tribes as an entheogen. In the Andean region of South America, the San Pedro cactus (Trichocereus macrogonus var. pachanoi, syn. Echinopsis pachanoi) has a long history of use, possibly as a traditional medicine. Archaeological studies have found evidence of use going back two thousand years, to Moche culture, Nazca culture, and Chavín culture. Although authorities of the Roman Catholic church attempted to suppress its use after the Spanish conquest, this failed, as shown by the Christian element in the common name "San Pedro cactus" – Saint Peter cactus. The name has its origin in the belief that just as St Peter holds the keys to heaven, the effects of the cactus allow users "to reach heaven while still on earth." In 2022, the Peruvian Ministry of Culture declared the traditional use of San Pedro cactus in northern Peru as cultural heritage.
Although people of Western culture have tended to use psychedelics for either psychotherapeutic or recreational reasons, most indigenous cultures, particularly in South America, have seemingly tended to use psychedelics for more supernatural reasons such as divination. This can often be related to "healing" or health as well but typically in the context of finding out what is wrong with the individual, such as using psychedelic states to "identify" a disease and/or its cause, locate lost objects, and identify a victim or even perpetrator of sorcery. In some cultures and regions, even psychedelics themselves, such as ayahuasca and the psychedelic lichen of eastern Ecuador (Dictyonema huaorani) that supposedly contains both 5-MeO-DMT and psilocybin, have also been used by witches and sorcerers to conduct their malicious magic, similarly to nightshade deliriants like brugmansia and latua.
Psychedelic therapy
Psychedelic therapy (or psychedelic-assisted therapy) is the proposed use of psychedelic drugs to treat mental disorders. As of 2021, psychedelic drugs are controlled substances in most countries and psychedelic therapy is not legally available outside clinical trials, with some exceptions.
The procedure for psychedelic therapy differs from that of therapies using conventional psychiatric medications. While conventional medications are usually taken without supervision at least once daily, in contemporary psychedelic therapy the drug is administered in a single session (or sometimes up to three sessions) in a therapeutic context. The therapeutic team prepares the patient for the experience beforehand and helps them integrate insights from the drug experience afterwards. After ingesting the drug, the patient normally wears eyeshades and listens to music to facilitate focus on the psychedelic experience, with the therapeutic team interrupting only to provide reassurance if adverse effects such as anxiety or disorientation arise.
As of 2022, the body of high-quality evidence on psychedelic therapy remains relatively small and more, larger studies are needed to reliably show the effectiveness and safety of psychedelic therapy's various forms and applications. On the basis of favorable early results, ongoing research is examining proposed psychedelic therapies for conditions including major depressive disorder, and anxiety and depression linked to terminal illness. The United States Food and Drug Administration has granted "breakthrough therapy" status, which expedites the assessment of promising drug therapies for potential approval, to psilocybin therapy for treatment-resistant depression and major depressive disorder.
Recreational
Recreational use of psychedelics has been common since the psychedelic era of the mid-1960s and continues to play a role in various festivals and events, including Burning Man. A survey published in 2013 found that 13.4% of American adults had used a psychedelic.
A June 2024 report by the RAND Corporation suggests psilocybin mushrooms may be the most prevalent psychedelic drug among adults in the United States. The RAND national survey indicated that 3.1% of U.S. adults reported using psilocybin in the past year. Roughly 12% of respondents acknowledged lifetime use of psilocybin, while a similar percentage reported having used LSD at some point in their lives. MDMA, also known as ecstasy, showed a lower prevalence of use at 7.6%. Notably, less than 1% of U.S. adults reported using any psychedelic drugs within the past month.
Microdosing
Psychedelic microdosing is the practice of using sub-threshold doses (microdoses) of psychedelics in an attempt to improve creativity, boost physical energy level, emotional balance, increase performance on problems-solving tasks and to treat anxiety, depression and addiction. The practice of microdosing has become more widespread in the 21st century with more people claiming long-term benefits from the practice.
A 2022 study recognized signatures of psilocybin microdosing in natural language and concluded that low amount of psychedelics have potential for application, and ecological observation of microdosing schedules.
Pharmacology
Mechanism of action
Most serotonergic psychedelics act as non-selective agonists of serotonin receptors, including of the serotonin 5-HT2 receptors, but often also of other serotonin receptors, such as the serotonin 5-HT1 receptors. They are thought to mediate their hallucinogenic effects specifically by activation of serotonin 5-HT2A receptors. Psychedelics, such as the tryptamines psilocin, DMT, and 5-MeO-DMT, the phenethylamines mescaline, DOM, and 2C-B, and ergolines and lysergamides like LSD, all act as agonists of the serotonin 5-HT2A receptors. Some psychedelics, such as phenethylamines like DOM and 2C-B, show high selectivity for the serotonin 5-HT2 receptors over other serotonin receptors. There is a very strong correlation between 5-HT2A receptor affinity and human hallucinogenic potency. In addition, the intensity of hallucinogenic effects in humans is directly correlated with the level of serotonin 5-HT2A receptor occupancy as measured with positron emission tomography (PET) imaging. Serotonin 5-HT2A receptor blockade with drugs like the semi-selective ketanserin and the non-selective risperidone can abolish the hallucinogenic effects of psychedelics in humans. However, studies with more selective serotonin 5-HT2A receptor antagonists, like pimavanserin, are still needed.
In animals, potency for stimulus generalization to the psychedelic DOM in drug discrimination tests is strongly correlated with serotonin 5-HT2A receptor affinity. Non-selective serotonin 5-HT2A receptor antagonists, like ketanserin and pirenperone, and selective serotonin 5-HT2A receptor antagonists, like volinanserin (MDL-100907), abolish the stimulus generalization of psychedelics in drug discrimination tests. Conversely, serotonin 5-HT2B and 5-HT2C receptor antagonists are ineffective. The potencies of serotonin 5-HT2 receptor antagonists in blocking psychedelic substitution are strongly correlated with their serotonin 5-HT2A receptor affinities. Highly selective serotonin 5-HT2A receptor agonists have recently been developed and show stimulus generalization to psychedelics, whereas selective serotonin 5-HT2C receptor agonists do not do so. The head-twitch response (HTR) is induced by serotonergic psychedelics and is a behavioral proxy of psychedelic-like effects in animals. The HTR is invariably induced by serotonergic psychedelics, is blocked by selective serotonin 5-HT2A receptor antagonists, and is abolished in serotonin 5-HT2A receptor knockout mice. In addition, there is a strong correlation between hallucinogenic potency in humans and potency in the HTR assay. Moreover, the HTR paradigm is one of the only animal tests that can distinguish between hallucinogenic serotonin 5-HT2A receptor agonists and non-hallucinogenic serotonin 5-HT2A receptor agonists, such as lisuride. In accordance with the preceding animal and human findings, it has been said that the evidence that the serotonin 5-HT2A receptor mediates the hallucinogenic effects of serotonergic psychedelics is overwhelming.
The serotonin 5-HT2A receptor activates several downstream signaling pathways. These include the Gq, β-arrestin2, and other pathways. Activation of both the Gq and β-arrestin2 pathways have been implicated in mediating the hallucinogenic effects of serotonergic psychedelics. However, subsequently, activation of the Gq pathway and not β-arrestin2 has been implicated. Interestingly, Gq signaling appeared to mediate hallucinogenic-like effects, whereas β-arrestin2 mediated receptor downregulation and tachyphylaxis. The lack of psychedelic effects with non-hallucinogenic serotonin 5-HT2A receptor agonists may be due to partial agonism of the serotonin 5-HT2A receptor with efficacy insufficient to produce psychedelic effects or may be due to biased agonism of the serotonin 5-HT2A receptor. There appears to be a threshold level of Gq activation (in terms of intrinsic activity, with >70%) required for production of hallucinogenic effects. Full agonists and partial agonists above this threshold are psychedelic 5-HT2A receptor agonists, whereas partial agonists below this threshold, such as lisuride, 2-bromo-LSD, 6-fluoro-DET, 6-MeO-DMT, and Ariadne, are non-hallucinogenic 5-HT2A receptor agonists. In addition, biased agonists that activate β-arrestin2 signaling but not Gq signaling, such as ITI-1549, IHCH-7086, and 25N-N1-Nap, are non-hallucinogenic serotonin 5-HT2A receptor agonists.
The hallucinogenic effects of serotonergic psychedelics may be critically mediated by serotonin 5-HT2A receptor activation in the medial prefrontal cortex (mPFC). Layer V pyramidal neurons in this area are especially discussed. Activation of serotonin 5-HT2A receptors in the mPFC results in marked excitatory and inhibitory effects as well as increased release of glutamate and GABA. Direct injection of serotonin 5-HT2A receptor agonists into the mPFC produces the HTR. Drugs that suppress glutamatergic activity in the mPFC, including AMPA receptor antagonists, metabotropic glutamate mGlu2/3 receptor agonists, μ-opioid receptor agonists, and adenosine A1 receptor agonists, block or suppress many of the neurochemical and behavioral effects of serotonergic psychedelics, including the HTR. Metabotropic glutamate mGlu2 receptors are primarily expressed as presynaptic autoreceptors and have inhibitory effects on glutamate release. Serotonergic psychedelics have been found to produce frontal cortex hyperactivity in humans in PET and single-photon emission computed tomography (SPECT) imaging studies. The PFC projects to many other cortical and subcortical brain areas, such as the locus coeruleus, nucleus accumbens, and amygdala, among others, and activation of the PFC by serotonergic psychedelics may thereby indirectly modulate these areas. In addition to the PFC, there is moderate to high expression of serotonin 5-HT2A receptors in the primary visual cortex (V1), as well as expression of the serotonin 5-HT2A receptor in other visual areas, and activation of these receptors may contribute to or mediate the visual effects of serotonergic psychedelics. Serotonergic psychedelics also directly or indirectly modulate a variety of other brain areas, like the claustrum, and this may be involved in their effects as well.
Serotonin, as well as drugs that increase serotonin levels, like the serotonin precursor 5-hydroxytryptophan (5-HTP), serotonin reuptake inhibitors, and serotonin releasing agents, are non-hallucinogenic in humans despite increasing activation of serotonin 5-HT2A receptors. Serotonin is a hydrophilic molecule which cannot easily cross biological membranes without active transport, and the serotonin 5-HT2A receptor is usually expressed as a cell surface receptor that is readily accessible to extracellular serotonin. The HTR, a behavioral proxy of psychedelic-like effects, appears to be mediated by activation of intracellularly expressed serotonin 5-HT2A receptors in a population of mPFC neurons that do not also express the serotonin transporter (SERT) and hence cannot be activated by serotonin. In contrast to serotonin, serotonergic psychedelics are more lipophilic than serotonin and are able to readily enter these neurons and activate the serotonin 5-HT2A receptors within them. Artificial expression of the SERT in this population of neurons in animals resulted in a serotonin releasing agent that doesn't normally produce the HTR being able to do so. Although serotonin itself is non-hallucinogenic, at very high concentrations achieved pharmacologically (e.g., injected into the brain or with massive doses of 5-HTP) it can produce psychedelic-like effects in animals by being metabolized by indolethylamine N-methyltransferase (INMT) into more lipophilic N-methylated tryptamines like N-methylserotonin and bufotenin (N,N-dimethylserotonin).
In addition to their hallucinogenic effects, serotonergic psychedelics may also produce a variety of other effects, including psychoplastogenic (i.e., neuroplasticity-enhancing), antidepressant, anxiolytic, empathy-enhancing or prosocial effects, anti-obsessional, anti-addictive, anti-inflammatory and immunomodulatory effects, analgesic effects, and/or antimigraine effects. While psychedelics themselves are also being clinically evaluated for these potential therapeutic benefits, non-hallucinogenic serotonin 5-HT2A receptor agonists, which are often analogues of serotonergic psychedelics, have been developed and are being studied for potential use in medicine in an attempt to provide some such benefits without hallucinogenic effects.
Although the hallucinogenic effects of serotonergic psychedelics are thought to be mediated by serotonin 5-HT2A receptor activation, interactions with other receptors, such as the serotonin 5-HT1A, 5-HT1B, 5-HT2B, and 5-HT2C receptors among many others, may additionally contribute to and modulate their effects. Interestingly, some psychedelics, such as LSD and psilocybin, have been claimed to act as positive allosteric modulators of the tropomyosin receptor kinase B (TrkB), one of the signaling receptors of brain-derived neurotrophic factor (BDNF). However, despite this apparent TrkB potentiation, the psychoplastogenic effects of serotonergic psychedelics, including dendritogenesis, spinogenesis, and synaptogenesis, appear to be mediated by activation of serotonin 5-HT2A receptors, whereas psychedelics do not generally stimulate neurogenesis.
Chemical families
The three major chemical groups of serotonergic psychedelics include the tryptamines, phenethylamines, and lysergamides, which each have different profiles of pharmacological activity.
Tryptamines
Tryptamines are derivatives of tryptamine and are structurally related to the monoamine neurotransmitter serotonin (also known as 5-hydroxytryptamine or 5-HT). Many tryptamines act as non-selective serotonin receptor agonists, including of the serotonin 5-HT2A receptor. Some tryptamines also act as monoamine releasing agents, including of serotonin, norepinephrine, and/or dopamine. Examples of psychedelic tryptamines include tryptamine, dimethyltryptamine (DMT), 5-MeO-DMT, psilocin, psilocybin, bufotenin, 5-MeO-DiPT, 5-MeO-MiPT, α-methyltryptamine (αMT), and 5-MeO-αMT, among many others.
Phenethylamines
Phenethylamines, as well as amphetamines (α-methylphenethylamines), are derivatives of β-phenethylamine and are structurally related to the monoamine neurotransmitters dopamine, norepinephrine, and epinephrine. Some phenethylamines and amphetamines, particularly those with methoxy and other substitions on the phenyl ring, are potent serotonin 5-HT2 receptor agonists, including of the serotonin 5-HT2A receptor, and can produce psychedelic effects. In contrast to phenethylamines and amphetamines generally, most psychedelic phenethylamines are not monoamine releasing agents. Examples of psychedelic phenethylamines and amphetamines include mescaline, the 2C drugs like 2C-B, 2C-D, 2C-E, and 2C-I, the DOx drugs like DOB, DOI, and DOM, certain MDxx drugs like MDA and MDMA (weak psychedelics), and the NBOMe (25x-NBx) drugs like 25C-NBOMe and 25I-NBOMe.
Lysergamides
Lysergamides are ergoline derivatives related to the ergot alkaloids. They are notable in containing both tryptamine and phenethylamine within their chemical structures. As such, ergolines and lysergamides may be considered structurally related to the monoamine neurotransmitters. Many ergolines and lysergamides act as highly promiscuous ligands of monoamine receptors, including of serotonin, dopamine, and adrenergic receptors. Some lysergamides are efficacious serotonin 5-HT2A receptor agonists and thereby produce psychedelic effects. Examples of psychedelic lysergamides include lysergic acid diethylamide (LSD), lysergic acid amide (LSA), 1P-LSD, ALD-52, ergonovine (ergometrine), and methylergometrine (methylergonovine), among others.
Psychedelic experiences
Although several attempts have been made, starting in the 19th and 20th centuries, to define common phenomenological structures of the effects produced by classic psychedelics, a universally accepted taxonomy does not yet exist. At lower doses, features of psychedelic experiences include sensory alterations, such as the warping of surfaces, shape suggestibility, pareidolia and color variations. Users often report intense colors that they have not previously experienced, and repetitive geometric shapes or form constants are common as well. Higher doses often cause intense and fundamental alterations of sensory (notably visual) perception, such as synesthesia or the experience of additional spatial or temporal dimensions. Tryptamines are well documented to cause classic psychedelic states, such as increased empathy, visual distortions (drifting, morphing, breathing, melting of various surfaces and objects), auditory hallucinations, ego dissolution or ego death with high enough dose, mystical, transpersonal and spiritual experiences, autonomous "entity" encounters, time distortion, closed eye hallucinations and complete detachment from reality with a high enough dose. Luis Luna describes psychedelic experiences as having a distinctly gnosis-like quality, and says that they offer "learning experiences that elevate consciousness and can make a profound contribution to personal development." Czech psychiatrist Stanislav Grof studied the effects of psychedelics like LSD early in his career and said of the experience, that it commonly includes "complex revelatory insights into the nature of existence… typically accompanied by a sense of certainty that this knowledge is ultimately more relevant and 'real' than the perceptions and beliefs we share in everyday life." Traditionally, the standard model for the subjective phenomenological effects of psychedelics has typically been based on LSD, with anything that is considered "psychedelic" evidently being compared to it and its specific effects.
During a speech on his 100th birthday, the inventor of LSD, Albert Hofmann said of the drug: "It gave me an inner joy, an open mindedness, a gratefulness, open eyes and an internal sensitivity for the miracles of creation... I think that in human evolution it has never been as necessary to have this substance LSD. It is just a tool to turn us into what we are supposed to be." With certain psychedelics and experiences, a user may also experience an "afterglow" of improved mood or perceived mental state for days or even weeks after ingestion in some cases. In 1898, the English writer and intellectual Havelock Ellis reported a heightened perceptual sensitivity to "the more delicate phenomena of light and shade and color" for a prolonged period of time after his exposure to mescaline. Good trips are reportedly deeply pleasurable, and typically involve intense joy or euphoria, a greater appreciation for life, reduced anxiety, a sense of spiritual enlightenment, and a sense of belonging or interconnectedness with the universe. Negative experiences, colloquially known as "bad trips," evoke an array of dark emotions, such as irrational fear, anxiety, panic, paranoia, dread, distrustfulness, hopelessness, and even suicidal ideation. While it is impossible to predict when a bad trip will occur, one's mood, surroundings, sleep, hydration, social setting, and other factors can be controlled (colloquially referred to as "set and setting") to minimize the risk of a bad trip. The concept of "set and setting" also generally appears to be more applicable to psychedelics than to other types of hallucinogens such as deliriants, hypnotics and dissociative anesthetics.
Classic psychedelics are considered to be those found in nature like psilocybin, DMT, mescaline, and LSD which is derived from naturally occurring ergotamine, and non-classic psychedelics are considered to be newer analogs and derivatives of pharmacophore lysergamides, tryptamine, and phenethylamine structures like 2C-B. Many of these psychedelics cause remarkably similar effects, despite their different chemical structure. However, many users report that the three major families have subjectively different qualities in the "feel" of the experience, which are difficult to describe. Some compounds, such as 2C-B, have extremely tight "dose curves", meaning the difference in dose between a non-event and an overwhelming disconnection from reality can be very slight. There can also be very substantial differences between the drugs; for instance, 5-MeO-DMT rarely produces the visual effects typical of other psychedelics.
Potential adverse effects
Despite the contrary perception of much of the public, psychedelic drugs are not addictive and are physiologically safe. As of 2016, there have been no known deaths due to overdose of LSD, psilocybin, or mescaline.
Risks do exist during an unsupervised psychedelic experience, however; Ira Byock wrote in 2018 in the Journal of Palliative Medicine that psilocybin is safe when administered to a properly screened patient and supervised by a qualified professional with appropriate set and setting. However, he called for an "abundance of caution" because in the absence of these conditions a range of negative reactions is possible, including "fear, a prolonged sense of dread, or full panic." He notes that driving or even walking in public can be dangerous during a psychedelic experience because of impaired hand-eye coordination and fine motor control. In some cases, individuals taking psychedelics have performed dangerous or fatal acts because they believed they possessed superhuman powers.
Psilocybin-induced states of mind share features with states experienced in psychosis, and while a causal relationship between psilocybin and the onset of psychosis has not been established as of 2011, researchers have called for investigation of the relationship. Many of the persistent negative perceptions of psychological risks are unsupported by the currently available scientific evidence, with the majority of reported adverse effects not being observed in a regulated and/or medical context. A population study on associations between psychedelic use and mental illness published in 2013 found no evidence that psychedelic use was associated with increased prevalence of any mental illness. In any case, induction of psychosis has been associated with psychedelics in small percentages of individuals, and the rates appear to be higher in people with schizophrenia.
Using psychedelics poses certain risks of re-experiencing of the drug's effects, including flashbacks and hallucinogen persisting perception disorder (HPPD). These non-psychotic effects are poorly studied, but the permanent symptoms (also called "endless trip") are considered to be rare.
Serotonin syndrome can be caused by combining psychedelics with other serotonergic drugs, including certain antidepressants, opioids, CNS stimulants (e.g. MDMA), 5-HT1 agonists (e.g. triptans), herbs and others.
Serotonergic psychedelics are agonists not only of the serotonin 5-HT2A receptor but also of the serotonin 5-HT2B receptor and other serotonin receptors. A potential risk of frequent repeated use of serotonergic psychedelics is cardiac fibrosis and valvulopathy caused by 5-HT2B receptor activation. However, single high doses or widely spaced doses (e.g., months) are widely thought to be safe and concerns about cardiac toxicity apply more to chronic psychedelic microdosing or very frequent use (e.g., weekly). Selective 5-HT2A receptor agonists that do not activate the 5-HT2B receptor or other serotonin receptors, such as 25CN-NBOH, DMBMPP, and LPH-5, have been developed and are being studied. Selective 5-HT2A receptor agonists are expected to avoid the cardiac risks of 5-HT2B receptor activation.
Potential therapeutic effects
Psychedelic substances which may have therapeutic uses include psilocybin, LSD, and mescaline. During the 1950s and 1960s, lack of informed consent in some scientific trials on psychedelics led to significant, long-lasting harm to some participants. Since then, research regarding the effectiveness of psychedelic therapy has been conducted under strict ethical guidelines, with fully informed consent and a pre-screening to avoid people with psychosis taking part. Although the history behind these substances has hindered research into their potential medicinal value, scientists are now able to conduct studies and renew research that was halted in the 1970s. Some research has shown that these substances have helped people with such mental disorders as obsessive-compulsive disorder (OCD), post-traumatic stress disorder (PTSD), alcoholism, depression, and cluster headaches.
It has long been known that psychedelics promote neurite growth and neuroplasticity and are potent psychoplastogens. There is evidence that psychedelics induce molecular and cellular adaptations related to neuroplasticity and that these could potentially underlie therapeutic benefits. Psychedelics have also been shown to have potent anti-inflammatory activity and therapeutic effects in animal models of inflammatory diseases including asthma, and cardiovascular disease and diabetes.
Surrounding culture
Psychedelic culture includes manifestations such as psychedelic music, psychedelic art, psychedelic literature, psychedelic film, and psychedelic festivals. Examples of psychedelic music would be rock bands like the Grateful Dead, Jefferson Airplane and The Beatles. Many psychedelic bands and elements of the psychedelic subculture originated in San Francisco during the mid to late 1960s.
Legal status
Many psychedelics are classified under Schedule I of the United Nations Convention on Psychotropic Substances of 1971 as drugs with the greatest potential to cause harm and no acceptable medical uses. In addition, many countries have analogue laws; for example, in the United States, the Federal Analogue Act of 1986 automatically forbids any drugs sharing similar chemical structures or chemical formulas to prohibited substances if sold for human consumption.
In July 2022, though, under the United States Food and Drug Administration, the drug psilocybin was on track to be approved of as a treatment for depression, and MDMA as a treatment for post-traumatic stress disorder.
U.S. states such as Oregon and Colorado have also instituted decriminalization and legalization measures for accessing psychedelics and states like New Hampshire are attempting to do the same. J.D. Tuccille argues that increasing rates of use of psychedelics in defiance of the law are likely to result in more widespread legalization and decriminalization of access to the substances in the United States (as has happened with alcohol and cannabis).
| Biology and health sciences | Recreational drugs | Health |
53332 | https://en.wikipedia.org/wiki/Saffron | Saffron | Saffron () is a spice derived from the flower of Crocus sativus, commonly known as the "saffron crocus". The vivid crimson stigma and styles, called threads, are collected and dried for use mainly as a seasoning and colouring agent in food. The saffron crocus was slowly propagated throughout much of Eurasia and was later brought to parts of North Africa, North America, and Oceania.
Saffron's taste and iodoform-like or hay-like fragrance result from the phytochemicals picrocrocin and safranal. It also contains a carotenoid pigment, crocin, which imparts a rich golden-yellow hue to dishes and textiles. Its recorded history is attested in a 7th-century BC Assyrian botanical treatise, and it has been traded and used for thousands of years. As of 2024, Iran produced some 90% of the world total for saffron. At US$5,000 per kg or higher, saffron has long been the world's costliest spice by weight.
Etymology
A degree of uncertainty surrounds the origin of the English word "saffron". It might stem from the 12th-century Old French term safran, which comes from the Latin word , from the Persian (, za'farān), from the Persian word zarparān () meaning "gold strung" (implying either the golden stamens of the flower or the golden colour it creates when used as flavour).
Species
Description
The domesticated saffron crocus, Crocus sativus, is an autumn-flowering perennial plant unknown in the wild. It probably descends from the eastern Mediterranean autumn-flowering Crocus cartwrightianus which is also known as "wild saffron" and is native to mainland Greece, Euboea, Crete, Skyros and some islands of the Cyclades. The similar species C. thomasii and C. pallasii were considered as other possible ancestors. As a genetically monomorphic clone incapable of seed production, it was slowly propagated by humans throughout much of Eurasia. Various origins had been suggested for saffron, including Iran, Greece, Mesopotamia. and Kashmir.
It is a sterile triploid form, which means that three homologous sets of chromosomes make up each specimen's genetic complement; C. sativus bears eight chromosomal bodies per set, making for 24 in total. Being sterile, the purple flowers of C. sativus fail to produce viable seeds; reproduction hinges on human assistance: clusters of corms, underground, bulb-like, starch-storing organs, must be dug up, divided, and replanted. A corm survives for one season, producing via vegetative division up to ten "cormlets" that can grow into new plants in the next season. The compact corms are small, brown globules that can measure as large as in diameter, have a flat base, and are shrouded in a dense mat of parallel fibres; this coat is referred to as the "corm tunic". Corms also bear vertical fibres, thin and net-like, that grow up to above the plant's neck.
The plant sprouts 5–11 white and non-photosynthetic leaves known as cataphylls. These membrane-like structures cover and protect 5 to 11 true leaves as they bud and develop on the crocus flower. The latter are thin, straight, and blade-like green foliage leaves, which are , in diameter, which either expand after the flowers have opened ("hysteranthous") or do so simultaneously with their blooming ("synanthous"). C. sativus cataphylls are suspected by some to manifest prior to blooming when the plant is irrigated relatively early in the growing season. Its floral axes, or flower-bearing structures, bear bracteoles, or specialised leaves, that sprout from the flower stems; the latter are known as pedicels. After aestivating in spring, the plant sends up its true leaves, each up to in length. Only in October, after most other flowering plants have released their seeds, do its brilliantly hued flowers develop; they range from a light pastel shade of lilac to a darker and more striated mauve. The flowers possess a sweet, honey-like fragrance. Upon flowering, the plants are in height and bear up to four flowers. A three-pronged style in length, emerges from each flower. Each prong terminates with a vivid crimson stigma, which is the distal end of a carpel.
Cultivation
The saffron crocus, unknown in the wild, probably descends from Crocus cartwrightianus. It is a triploid that is "self-incompatible" and male sterile; it undergoes aberrant meiosis and is hence incapable of independent sexual reproduction—all propagation is by vegetative multiplication via manual "divide-and-set" of a starter clone or by interspecific hybridisation.
Crocus sativus thrives in the Mediterranean maquis, an ecotype superficially resembling the North American chaparral, and similar climates where hot and dry summer breezes sweep semi-arid lands. It can nonetheless survive cold winters, tolerating frosts as low as and short periods of snow cover. Some reports suggest saffron can tolerate an air temperature range from −22 to 40 °C. Irrigation is required if grown outside of moist environments such as Kashmir, where annual rainfall averages ; saffron-growing regions in Greece ( annually) and Spain () are far drier than the main cultivating Iranian regions. What makes this possible is the timing of the local wet seasons; generous spring rains and drier summers are optimal. Rain immediately preceding flowering boosts saffron yields; rainy or cold weather during flowering promotes disease and reduces yields. Persistently damp and hot conditions harm the crops, and rabbits, rats, and birds cause damage by digging up corms. Nematodes, leaf rusts, and corm rot pose other threats. Yet Bacillus subtilis inoculation may provide some benefit to growers by speeding corm growth and increasing stigma biomass yield.
The plants fare poorly in shady conditions; they grow best in full sunlight. Fields that slope towards the sunlight are optimal (i.e., south-sloping in the Northern Hemisphere). Planting is mostly done in June in the Northern Hemisphere, where corms are lodged deep; its roots, stems, and leaves can develop between October and February. Planting depth and corm spacing, in concert with climate, are critical factors in determining yields. Mother corms planted deeper yield higher-quality saffron, though they form fewer flower buds and daughter corms. Italian growers optimise thread yield by planting deep and in rows apart; depths of optimise flower and corm production. Greek, Moroccan, and Spanish growers employ distinct depths and spacings that suit their locales.
C. sativus prefers friable, loose, low-density, well-watered, and well-drained clay-calcareous soils with high organic content. Traditional raised beds promote good drainage. Soil organic content was historically boosted via application of some of manure. Afterwards, and with no further manure application, corms were planted. After a period of dormancy through the summer, the corms send up their narrow leaves and begin to bud in early autumn. Only in mid-autumn do they flower. Harvests are by necessity a speedy affair: after blossoming at dawn, flowers quickly wilt as the day passes. All plants bloom within a window of one or two weeks. Stigmas are dried quickly upon extraction and (preferably) sealed in airtight containers.
Harvesting
The high retail value of saffron is maintained on world markets because of labour-intensive harvesting methods, which require some – equivalently, . Forty hours of labour are needed to pick 150,000 flowers.
One freshly picked crocus flower yields on average 30 mg of fresh saffron or 7 mg dried; roughly 150 flowers yield of dry saffron threads; to produce of dried saffron, of flowers are needed; the yield of dried spice from fresh saffron is only .
Spice
Phytochemistry and sensory properties
Saffron contains some 28 volatile and aroma-yielding compounds, dominated by ketones and aldehydes. Its main aroma-active compounds are safranal – the main compound responsible for saffron aroma – 4-ketoisophorone, and dihydrooxophorone. Saffron also contains nonvolatile phytochemicals, including the carotenoids zeaxanthin, lycopene, various α- and β-carotenes, as well as crocetin and its glycoside crocein, which are the most biologically active components. Because crocetin is smaller and more water-soluble than the other carotenoids, it is more rapidly absorbed.
The yellow-orange colour of saffron is primarily the result of α-crocin. This crocin is trans-crocetin di-(β-D-gentiobiosyl) ester; it bears the systematic (IUPAC) name 8,8-diapo-8,8-carotenoic acid. This means that the crocin underlying saffron's aroma is a digentiobiose ester of the carotenoid crocetin. Crocins themselves are a series of hydrophilic carotenoids that are either monoglycosyl or diglycosyl polyene esters of crocetin. Crocetin is a conjugated polyene dicarboxylic acid that is hydrophobic, and thus oil-soluble. When crocetin is esterified with two water-soluble gentiobioses, which are sugars, a product results that is itself water-soluble. The resultant α-crocin is a carotenoid pigment that may make up more than 10% of dry saffron's mass. The two esterified gentiobioses make α-crocin ideal for colouring water-based and non-fatty foods such as rice dishes.
The bitter glucoside picrocrocin is responsible for saffron's pungent flavour. Picrocrocin (chemical formula: ; systematic name: 4-(β-D-glucopyranosyloxy)-2,6,6-trimethylcyclohex-1-ene-1-carbaldehyde) is a union of an aldehyde sub-molecule known as safranal (systematic name: 2,6,6-trimethylcyclohexa-1,3-diene-1-carbaldehyde) and a carbohydrate. It has insecticidal and pesticidal properties, and may comprise up to 4% of dry saffron. Picrocrocin is a truncated version of the carotenoid zeaxanthin that is produced via oxidative cleavage, and is the glycoside of the terpene aldehyde safranal.
When saffron is dried after its harvest, the heat, combined with enzymatic action, splits picrocrocin to yield D–glucose and a free safranal molecule. Safranal, a volatile oil, gives saffron much of its distinctive aroma. Safranal is less bitter than picrocrocin and may comprise up to 70% of dry saffron's volatile fraction in some samples. A second molecule underlying saffron's aroma is 2-hydroxy-4,4,6-trimethyl-2,5-cyclohexadien-1-one, which produces a scent described as saffron, dried hay-like. Chemists find this is the most powerful contributor to saffron's fragrance, despite its presence in a lesser quantity than safranal. Dry saffron is highly sensitive to fluctuating pH levels, and rapidly breaks down chemically in the presence of light and oxidising agents. It must, therefore, be stored in air-tight containers to minimise contact with atmospheric oxygen. Saffron is somewhat more resistant to heat.
Grades and ISO 3632 categories
Saffron is not all of the same quality and strength. Strength is related to several factors including the amount of style picked along with the red stigma. Age of the saffron is also a factor. More style included means the saffron is less strong gram for gram because the colour and flavour are concentrated in the red stigmas. Saffron from Iran, Spain, and Kashmir is classified into various grades according to the relative amounts of red stigma and yellow styles it contains. Grades of Iranian saffron are: sargol (, red stigma tips only, strongest grade), pushal or pushali (red stigmas plus some yellow style, lower strength), "bunch" saffron (red stigmas plus large amount of yellow style, presented in a tiny bundle like a miniature wheatsheaf) and konge (yellow style only, claimed to have aroma but with very little, if any, colouring potential). Grades of Spanish saffron are coupé (the strongest grade, like Iranian sargol), mancha (like Iranian pushal), and in order of further decreasing strength rio, standard and sierra saffron. The word mancha in the Spanish classification can have two meanings: a general grade of saffron or a very high quality Spanish-grown saffron from a specific geographical origin. Real Spanish-grown La Mancha saffron has PDO protected status and this is displayed on the product packaging. Spanish growers fought hard for Protected Status because they felt that imports of Iranian saffron re-packaged in Spain and sold as "Spanish Mancha saffron" were undermining the genuine La Mancha brand. Similar was the case in Kashmir where imported Iranian saffron is mixed with local saffron and sold as "Kashmir brand" at a higher price. In Kashmir, saffron is mostly classified into two main categories called mongra (stigma alone) and lachha (stigmas attached with parts of the style). Countries producing less saffron do not have specialised words for different grades and may only produce one grade. Artisan producers in Europe and New Zealand have offset their higher labour charges for saffron harvesting by targeting quality, only offering extremely high-grade saffron.
In addition to descriptions based on how the saffron is picked, saffron may be categorised under the international standard ISO 3632 after laboratory measurement of crocin (responsible for saffron's colour), picrocrocin (taste), and safranal (fragrance or aroma) content. However, often there is no clear grading information on the product packaging and little of the saffron readily available in the UK is labelled with ISO category. This lack of information makes it hard for customers to make informed choices when comparing prices and buying saffron.
Under ISO 3632, determination of non-stigma content ("floral waste content") and other extraneous matter such as inorganic material ("ash") are also key. Grading standards are set by the International Organization for Standardization, a federation of national standards bodies. ISO 3632 deals exclusively with saffron and establishes three categories: III (poorest quality), II, and I (finest quality). Formerly there was also category IV, which was below category III. Samples are assigned categories by gauging the spice's crocin and picrocrocin content, revealed by measurements of specific spectrophotometric absorbance. Safranal is treated slightly differently and rather than there being threshold levels for each category, samples must give a reading of 20–50 for all categories.
These data are measured through spectrophotometry reports at certified testing laboratories worldwide. Higher absorbances imply greater levels of crocin, picrocrocin and safranal, and thus a greater colouring potential and therefore strength per gram. The absorbance reading of crocin is known as the "colouring strength" of that saffron. Saffron's colouring strength can range from lower than 80 (for all category IV saffron) up to 200 or greater (for category I). The world's finest samples (the selected, most red-maroon, tips of stigmas picked from the finest flowers) receive colouring strengths in excess of 250, making such saffron over three times more powerful than category IV saffron. Market prices for saffron types follow directly from these ISO categories. Sargol and coupé saffron would typically fall into ISO 3632 category I. Pushal and Mancha would probably be assigned to category II. On many saffron packaging labels, neither the ISO 3632 category nor the colouring strength (the measurement of crocin content) is displayed.
However, many growers, traders, and consumers reject such lab test numbers. Some people prefer a more holistic method of sampling batches of threads for taste, aroma, pliability, and other traits in a fashion similar to that practised by experienced wine tasters.
Adulteration
Despite attempts at quality control and standardisation, an extensive history of saffron adulteration, particularly among the cheapest grades, continues into modern times. Adulteration was first documented in Europe's Middle Ages, when those found selling adulterated saffron were executed under the Safranschou code. Typical methods include mixing in extraneous substances like beetroot, pomegranate fibres, red-dyed silk fibres, or the saffron crocus's tasteless and odourless yellow stamens. Other methods included dousing saffron fibres with viscid substances like honey or vegetable oil to increase their weight. Powdered saffron is more prone to adulteration, with turmeric, paprika, and other powders used as diluting fillers. Adulteration can also consist of selling mislabelled mixes of different saffron grades. Thus, high-grade Kashmiri saffron is often sold and mixed with cheaper Iranian imports; these mixes are then marketed as pure Kashmiri saffron. Safflower is a common substitute sometimes sold as saffron. The spice is reportedly counterfeited with horse hair, corn silk, or shredded paper. Tartrazine or sunset yellow have been used to colour counterfeit powdered saffron.
In recent years, saffron adulterated with the colouring extract of gardenia fruits has been detected in the European market. This form of fraud is difficult to detect due to the presence of flavonoids and crocines in the gardenia-extracts similar to those naturally occurring in saffron. Detection methods have been developed by using HPLC and mass spectrometry to determine the presence of geniposide, a compound present in the fruits of gardenia, but not in saffron.
Types
The various saffron crocus cultivars give rise to thread types that are often regionally distributed and characteristically distinct. Varieties (not varieties in the botanical sense) from Spain, including the tradenames "Spanish Superior" and "Creme", are generally mellower in colour, flavour, and aroma; they are graded by government-imposed standards. Italian varieties are slightly more potent than Spanish. Greek saffron produced in the town of Krokos is PDO protected due to its particularly high-quality colour and strong flavour. Various "boutique" crops are available from New Zealand, France, Switzerland, England, the United States, and other countries—some of them organically grown. In the US, Pennsylvania Dutch saffron—known for its "earthy" notes—is marketed in small quantities.
Consumers may regard certain cultivars as "premium" quality. The "Aquila" saffron, or zafferano dell'Aquila, is defined by high safranal and crocin content, distinctive thread shape, unusually pungent aroma, and intense colour; it is grown exclusively on eight hectares in the Navelli Valley of Italy's Abruzzo region, near L'Aquila. It was first introduced to Italy by a Dominican friar from inquisition-era Spain. But the biggest saffron cultivation in Italy is in San Gavino Monreale, Sardinia, where it is grown on 40 hectares, representing 60% of Italian production; it too has unusually high crocin, picrocrocin, and safranal content.
Another is the "Mongra" or "Lacha" saffron of Kashmir (Crocus sativus 'Cashmirianus'), which is among the most difficult for consumers to obtain. Repeated droughts, blights, and crop failures in Kashmir combined with an Indian export ban, contribute to its prohibitive overseas prices. Kashmiri saffron is recognizable by its dark maroon-purple hue, making it among the world's darkest. In 2020, Kashmir Valley saffron was certified with a geographical indication from the Government of India.
World production
Almost all saffron grows in a belt from Spain in the west to India in the east. Iran is responsible for around 88% of global production. In 2018, Iran cultivated an area of producing 174 tonnes from a productivity of 4 kg/ha. Afghanistan comes second, which produced over 67 tons in 2023.
Spain is the third largest producer, while the United Arab Emirates, Greece, the Indian subcontinent and Morocco are among minor producers.
According to the statistics for saffron trade in 2019, Iran was ranked as the world's largest producer of saffron, supplying 430 tons of the total 450 tons of saffron produced worldwide and is expected to reach 500 tons in 2020. India, producing only 22 tons of saffron annually, ranked second. Other countries reported based on their share in global saffron production included Greece (7.2 tons), Afghanistan (6 tons), Morocco (2.6 tons), Spain (2.3 tons), Italy (1 ton), China (1 ton), and Azerbaijan (0.23 ton).
Trade
Saffron prices at wholesale and retail rates range from . In Western countries, the average retail price in 1974 was . In February 2013, a retail bottle containing could be purchased for $16.26 or the equivalent of , or as little as about in larger quantities. There are between . Vivid crimson colouring, slight moistness, elasticity, and lack of broken-off thread debris are all traits of fresh saffron.
Uses
Saffron has a long history of use in traditional medicine. Saffron has also been used as a fabric dye, particularly in China and India, and in perfumery. It is used for religious purposes in India.
In the European E number categorisation for food elements and additives, Saffron is coded as E164.
Consumption
Saffron's aroma is often described by connoisseurs as reminiscent of metallic honey with grassy or hay-like notes, while its taste has also been noted as hay-like and sweet. Saffron also contributes a luminous yellow-orange colouring to foods. Saffron is widely used in Persian, Indian, European, and Arab cuisines. Confectioneries and liquors also often include saffron. Saffron is used in dishes ranging from the jewelled rice and khoresh of Iran, the Milanese risotto of Italy, the paella of Spain, the bouillabaisse of France, to the biryani with various meat accompaniments in South Asia. Saffron is also used in the preparation of the Golden Ham, a precious dry-cured ham made with saffron from San Gimignano. Common saffron substitutes include safflower (Carthamus tinctorius, which is often sold as "Portuguese saffron" or "açafrão"), annatto, and turmeric (Curcuma longa). In Medieval Europe, turmeric was also known as "Indian saffron" because of its yellow-orange colour.
Nutrition
Dried saffron is 65% carbohydrates, 6% fat, 11% protein (table) and 12% water. In one tablespoon (2 grams; a quantity much larger than is likely to be ingested in normal use) manganese is present as 29% of the Daily Value, while other micronutrients have negligible content (table).
Toxicity
Ingesting less than of saffron is not toxic for humans, but doses greater than can become increasingly toxic. Mild toxicity includes dizziness, nausea, vomiting, and diarrhea, whereas at higher doses there can be reduced platelet count and spontaneous bleeding.
Storage
Saffron will not spoil, but will lose flavour within six months if not stored in an airtight, cool and dark place. Freezer storage can maintain flavour for up to two years.
Research
As of 2020, saffron constituents, such as crocin, crocetin, and safranal, were under preliminary research for their potential to affect mental depression. Saffron has also been studied for its possible effect on cardiovascular risk factors, and in erectile dysfunction.
History
Saffron likely originated in Iran, Greece, Mesopotamia, or Kashmir. Harold McGee states that it was domesticated in or near Greece during the Bronze Age. C. sativus is probably a triploid form of Crocus cartwrightianus, which is also known as "wild saffron". Saffron crocus was slowly propagated by humans throughout much of Eurasia and was later brought to parts of North Africa, North America, and Oceania.
Several wild species of Crocus similar to the commercial plant are known to have been harvested in recent times for use as saffron. Crocus ancyrensis was used to make saffron in Sivas in Central Turkey, the corms were also eaten. Crocus cartwrightianus was harvested on Andros in the islands of the Cyclades, for medicinal purposes and the stigmas for making a pigment called Zafran. Crocus longiflorus stigmas were used for saffron in Sicily. Crocus thomasii stigmas were used to flavour dishes around Taranto, South Italy. In Syria the stigmas of an unknown wild species were collected by women and children, sun-dried and pressed into small tablets which were sold in the Bazaars. Not all ancient depictions or descriptions of saffron spice or flowers are certain to be the same species as the modern commercial species used for spice.
West Asia
Saffron was detailed in a 7th-century BC Assyrian botanical reference compiled under Ashurbanipal. Documentation of saffron's use over the span of 3,500 years has been uncovered. Saffron-based pigments have indeed been found in 50,000-year-old depictions of prehistoric places in northwest Iran. The Sumerians later used wild-growing saffron in their remedies and magical potions. Saffron was an article of long-distance trade before the Minoan palace culture's 2nd millennium BC peak. Ancient Persians cultivated Persian saffron (Crocus sativus var. haussknechtii now called Crocus haussknechtii by botanists) in Derbent, Isfahan, and Khorasan by the 10th century BC. At such sites, saffron threads were woven into textiles, ritually offered to divinities, and used in dyes, perfumes, medicines, and body washes. Saffron threads would thus be scattered across beds and mixed into hot teas as a curative for bouts of melancholy. Non-Persians also feared the Persians' usage of saffron as a drugging agent and aphrodisiac. During his Asian campaigns, Alexander the Great used Persian saffron in his infusions, rice, and baths as a curative for battle wounds. Alexander's troops imitated the practice from the Persians and brought saffron-bathing to Greece.
South Asia
Conflicting theories explain saffron's arrival in South Asia. Kashmiri and Chinese accounts date its arrival anywhere between 2500 and 900 years ago. Historians studying ancient Persian records date the arrival to sometime prior to 500 BC, attributing it to a Persian transplantation of saffron corms to stock new gardens and parks. Phoenicians then marketed Kashmiri saffron as a dye and a treatment for melancholy. Its use in foods and dyes subsequently spread throughout South Asia. Buddhist monks wear saffron-coloured robes; however, the robes are not dyed with costly saffron but turmeric, a less expensive dye, or jackfruit. Monks' robes are dyed the same colour to show equality with each other, and turmeric or ochre were the cheapest, most readily available dyes. Gamboge is also used to dye the robes.
East Asia
Some historians believe that saffron came to China with Mongol invaders from Persia. Yet it is mentioned in ancient Chinese medical texts, including the forty-volume Shennong Bencaojing, a pharmacopoeia written around 300–200 BC. Traditionally credited to the legendary Yan Emperor and the deity Shennong, it discusses 252 plant-based medical treatments for various disorders. Nevertheless, around the 3rd century AD, the Chinese were referring to it as having a Kashmiri provenance. According to the herbalist Wan Zhen, "the habitat of saffron is in Kashmir, where people grow it principally to offer it to the Buddha". Wan also reflected on how it was used in his time: "The flower withers after a few days, and then the saffron is obtained. It is valued for its uniform yellow colour. It can be used to aromatise wine."
South East Mediterranean
Minoan depictions of saffron are now considered to be Crocus cartwrightianus. The Minoans portrayed saffron in their palace frescoes by 1600–1500 BC; they hint at its possible use as a therapeutic drug. Ancient Greek legends told of sea voyages to Cilicia, where adventurers sought what they believed were the world's most valuable threads. Another legend tells of Crocus and Smilax, whereby Crocus is bewitched and transformed into the first saffron crocus. Ancient perfumers in Egypt, physicians in Gaza, townspeople in Rhodes, and the Greek hetaerae courtesans used saffron in their scented waters, perfumes and potpourris, mascaras and ointments, divine offerings, and medical treatments.
In late Ptolemaic Egypt, Cleopatra used saffron in her baths so that lovemaking would be more pleasurable. Egyptian healers used saffron as a treatment for all varieties of gastrointestinal ailments. Saffron was also used as a fabric dye in such Levantine cities as Sidon and Tyre in Lebanon. Aulus Cornelius Celsus prescribes saffron in medicines for wounds, cough, colic, and scabies, and in the mithridatium.
Western Europe
Saffron was a notable ingredient in certain Roman recipes such as jusselle and conditum. Such was the Romans' love of saffron that Roman colonists took it with them when they settled in southern Gaul, where it was extensively cultivated until Rome's fall. With this fall, European saffron cultivation plummeted. Competing theories state that saffron only returned to France with 8th-century AD Moors or with the Avignon papacy in the 14th century AD. Similarly, the spread of Islamic civilisation may have helped reintroduce the crop to Spain and Italy.
The 14th-century Black Death caused demand for saffron-based medicaments to peak, and Europe imported large quantities of threads via Venetian and Genoan ships from southern and Mediterranean lands such as Rhodes. The theft of one such shipment by noblemen sparked the fourteen-week-long Saffron War. The conflict and resulting fear of rampant saffron piracy spurred corm cultivation in Basel; it thereby grew prosperous. The crop then spread to Nuremberg, where endemic and insalubrious adulteration brought on the Safranschou code—whereby culprits were variously fined, imprisoned, and executed. Meanwhile, cultivation continued in southern France, Italy, and Spain.
Direct archaeological evidence of mediaeval saffron consumption in Scandinavia comes from the wreck of the royal Danish-Norwegian flagship, Gribshunden. The ship sank in 1495 while on a diplomatic mission to Sweden. Excavations in 2021 revealed concentrations of saffron threads and small "pucks" of compressed saffron powder, along with fresh ginger, cloves, and pepper. Surprisingly, the saffron retained its distinctive odour even after more than 500 years of submersion in the Baltic Sea.
The Essex town of Saffron Walden, named for its new specialty crop, emerged as a prime saffron growing and trading centre in the 16th and 17th centuries but cultivation there was abandoned; saffron was re-introduced around 2013 as well as other parts of the UK (Cheshire).
The Americas
Europeans introduced saffron to the Americas when immigrant members of the Schwenkfelder Church left Europe with a trunk containing its corms. Church members had grown it widely in Europe. By 1730, the Pennsylvania Dutch cultivated saffron throughout eastern Pennsylvania. Spanish colonies in the Caribbean bought large amounts of this new American saffron, and high demand ensured that saffron's list price on the Philadelphia commodities exchange was equal to gold. Trade with the Caribbean later collapsed in the aftermath of the War of 1812, when many saffron-bearing merchant vessels were destroyed. Yet the Pennsylvania Dutch continued to grow lesser amounts of saffron for local trade and use in their cakes, noodles, and chicken or trout dishes. American saffron cultivation survives into modern times, mainly in Lancaster County, Pennsylvania.
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| Biology and health sciences | Monocots | null |
53335 | https://en.wikipedia.org/wiki/Endometriosis | Endometriosis | Endometriosis is a disease in which cells like those in the endometrium, the layer of tissue that normally covers the inside of the uterus, grow outside the uterus. It occurs in humans and a limited number of menstruating mammals. Endometrial tissue most often grows on or around reproductive organs such as the ovaries and fallopian tubes, on the outside surface of the uterus, or on the tissues surrounding the uterus and the ovaries (peritoneum). It can also grow on other organs in the pelvic region like the bowels, stomach, bladder, or the cervix. Rarely, it can also occur in other parts of the body.
Symptoms can be very different from person to person, varying in range and intensity. About 25% of individuals have no symptoms, while for some it can be a debilitating disease. Common symptoms include pelvic pain, heavy and painful periods, pain with bowel movements, painful urination, pain during sexual intercourse and infertility. Nearly half of those affected have chronic pelvic pain, while 70% feel pain during menstruation. Up to half of affected individuals are infertile. Besides physical symptoms, endometriosis can have an effect on the mental health and social life of people.
Diagnosis is usually based on symptoms and medical imaging; however, a definitive diagnosis is made through laparoscopy (excision is the gold standard) and biopsy. Other causes of similar symptoms include pelvic inflammatory disease, irritable bowel syndrome, interstitial cystitis, and fibromyalgia. Endometriosis is often misdiagnosed and many patients report being incorrectly told their symptoms are trivial or normal. Patients with endometriosis see an average of seven physicians before receiving a correct diagnosis, with an average delay of 6.7 years between the onset of symptoms and surgically obtained biopsies, the gold standard for diagnosing the condition.
It is estimated that globally around 10% of the female population, 190 million women are affected by endometriosis. Ethnic differences have been observed in endometriosis, as Southeast Asian and East Asian women are significantly more likely than White women to be diagnosed with endometriosis.
The exact cause of endometriosis is not known. Possible causes include problems with menstrual period flow, genetic factors, hormones, and problems with the immune system. Endometriosis is associated with elevated levels of the female sex hormone estrogen, as well as estrogen receptor sensitivity. Estrogen exposure worsens the inflammatory symptoms of endometriosis by stimulating an immune response.
While there is no cure for endometriosis, several treatments may improve symptoms. This may include pain medication, hormonal treatments or surgery. The recommended pain medication is usually a non-steroidal anti-inflammatory drug (NSAID), such as naproxen. Taking the active component of the birth control pill continuously or using an intrauterine device with progestogen may also be useful. Gonadotropin-releasing hormone agonist (GnRH agonist) may improve the ability of those who are infertile to conceive. Surgical removal of endometriosis may be used to treat those whose symptoms are not manageable with other treatments. Surgeons use ablation or excision to remove endometriosis lesions. However, excision is the gold standard of treatment for endometriosis, as it involves cutting out the lesions, as opposed to ablation, which is the burning of the lesions, which leaves no samples for biopsy to confirm endometriosis.
Signs and symptoms
Pain and infertility are common symptoms, although 20–25% of affected women are asymptomatic. Presence of pain symptoms are associated with the type of endometrial lesions as 50% of women with typical lesions, 10% of women with cystic ovarian lesions, and 5% of women with deep endometriosis do not have pain.
Pelvic pain
A major symptom of endometriosis is recurring pelvic pain. The pain can range from mild to severe cramping or stabbing pain that occurs on both sides of the pelvis, in the lower back and rectal area, and even down the legs. The amount of pain a person feels correlates weakly with the extent or stage (1 through 4) of endometriosis, with some individuals having little or no pain despite having extensive endometriosis or endometriosis with scarring, while others may have severe pain even though they have only a few small areas of endometriosis. The most severe pain is typically associated with menstruation. Pain can also start a week before a menstrual period, during and even a week after a menstrual period, or it can be constant. The pain can be debilitating and result in emotional stress. Symptoms of endometriosis-related pain may include:
Dysmenorrhea (64%) – painful, sometimes disabling cramps during the menstrual period; pain may get worse over time (progressive pain), also lower back pains linked to the pelvis
Chronic pelvic pain – typically accompanied by lower back pain or abdominal pain
Dyspareunia – painful sexual intercourse
Dysuria – urinary urgency, frequency, and sometimes painful voiding
Mittelschmerz – pain associated with ovulation
Bodily movement pain – present during exercise, standing, or walking
Compared with patients with superficial endometriosis, those with deep disease appear to be more likely to report shooting rectal pain and a sense of their insides being pulled down. Individual pain areas and pain intensity appear to be unrelated to the surgical diagnosis, and the area of pain unrelated to the area of endometriosis.
There are multiple causes of pain. Endometriosis lesions react to hormonal stimulation and may "bleed" during menstruation. The blood accumulates locally if it is not cleared shortly by the immune, circulatory, and lymphatic systems. This accumulation can lead to swelling, which triggers inflammation with the activation of cytokines, resulting in pain. Another source of pain is organ dislocation that arises from adhesion binding internal organs together. The ovaries, the uterus, the oviducts, the peritoneum, and the bladder can all be bound together. Pain triggered in this way can last throughout the menstrual cycle, not just during menstrual periods.
Additionally, endometriotic lesions can develop their own nerve supply, creating a direct and two-way interaction between lesions and the central nervous system. This interaction can produce a variety of individual differences in pain that, in some cases, become independent of the disease itself. Nerve fibers and blood vessels are thought to grow into endometriosis lesions by a process known as neuroangiogenesis.
Infertility
About a third of women with infertility have endometriosis. Among those with endometriosis, about 40% are infertile. The pathogenesis of infertility varies by disease stage: in early-stage disease, it is hypothesised to result from an inflammatory response that impairs various aspects of conception, whereas in later stages, distorted pelvic anatomy and adhesions contribute to impaired fertilisation.
Other
Bowel endometriosis may include symptoms like diarrhea, constipation, tenesmus, dyschezia, and, rarely, rectal bleeding. Other symptoms include chronic fatigue, nausea and vomiting, migraines, low-grade fevers, heavy (44%) and/or irregular periods (60%), and hypoglycemia. Endometriosis is associated with certain types of cancers, notably some types of ovarian cancer, non-Hodgkin's lymphoma and brain cancer. Endometriosis is however unrelated to endometrial cancer.
Rarely, endometriosis can cause endometrium-like tissue to be found in other parts of the body. Thoracic endometriosis occurs when endometrium-like tissue implants in the lungs or pleura. Manifestations of this include coughing up blood, a collapsed lung, or bleeding into the pleural space. Endometriosis may also affect the nearby colon which in rare situations may progress to partial obstruction requiring emergency surgery.
Stress may be a contributing factor or a consequence of endometriosis.
Complications
Physical health
Complications of endometriosis include internal scarring, adhesions, pelvic cysts, ovarian chocolate cysts, ruptured cysts, and bowel and ureter obstruction resulting from pelvic adhesions. Endometriosis-associated infertility may result from scar formation and anatomical distortions caused by the condition.
Ovarian endometriosis may complicate pregnancy through decidualization, abscess formation and/or rupture.
Thoracic endometriosis can be associated with recurrent thoracic endometriosis syndrome which manifests during menstrual periods. It includes catamenial pneumothorax in 73% of women, catamenial hemothorax in 14%, catamenial hemoptysis in 7%, and pulmonary nodules in 6%.
A 20-year study involving 12,000 women with endometriosis found that individuals under 40 are three times more likely to develop heart problems compared to their healthy peers.
A study indicated that 39% of women with surgically confirmed non-graded endometriosis had a 270% higher risk for ectopic pregnancy and a 76% higher risk for miscarriage compared to their peers. For women with deep endometriosis (>5 mm invasion, ASRM Stage II and higher), the risk of miscarriage increased by 298%.
Women with endometriosis also face a significantly increased risk of experiencing ante- and postpartum hemorrhage as well as a 170% increased risk of severe pre-eclampsia during pregnancy.
Endometriosis slightly increases the risk (about 1% or less) of developing ovarian, breast and thyroid cancers compared to women without the condition.
The mortality rates associated with endometriosis are low, with unadjusted and age-standardized death rates of 0.1 and 0.0 per 100,000, respectively.
Sciatic endometriosis also called catamenial or cyclical sciatica, is a rare form where endometriosis affects the sciatic nerve. Diagnosis is usually confirmed through MRI or CT-myelography.
Endometriosis can also impact a woman's fetus or neonate, increasing the risks for congenital malformations, preterm delivery and higher neonatal death rates.
Endometriosis can lead to ovarian cysts (endometriomas), adhesions, and damage to the fallopian tubes or ovaries, all of which can interfere with ovulation and fertilization. Treatment for endometriosis often includes hormonal therapies, pain management, and in some cases, surgery to remove the endometrial tissue. For women who struggle with infertility due to endometriosis, assisted reproductive technologies such as in vitro fertilization (IVF) may be recommended, sometimes in combination with surgical treatment to improve fertility outcomes.
Mental health
"Endometriosis is associated with an elevated risk of developing depression and anxiety disorders". Studies suggest this is partially due to the pelvic pain experienced by endometriosis patients.
Mental health concerns like depression and anxiety can also result due to poor diagnostic procedures related to cultural norms where women's concerns are devalued or ignored, especially by medical professionals.
Risk factors
Genetics
Endometriosis is a heritable condition influenced by both genetic and environmental factors, a genetic disorder of polygenic/multifactorial inheritance acquired via affected genes from either a person's father or mother. For example, children or siblings of women with endometriosis are at higher risk of developing endometriosis themselves; low progesterone levels may be genetic, and may contribute to a hormone imbalance. Individuals with an affected first-degree relative have an approximate six-fold increase incidence of endometriosis.
Inheritance is significant, but not the sole risk factor for endometriosis. Studies attribute 50% of risk to genetics, the other 50% likely to environmental factors. It has been proposed that endometriosis may result from a series of multiple mutations, within target genes, in a mechanism similar to the development of cancer. In this case, the mutations may be either somatic or heritable.
A 2019 genome-wide association study (GWAS) review enumerated 36 genes with mutations associated with endometriosis development. Nine chromosome loci were robustly replicated:
There are many findings of altered gene expression and epigenetics, but both of these can also be a secondary result of, for example, environmental factors and altered metabolism. Examples of altered gene expression include that of miRNAs.
Environmental toxins
Some factors associated with endometriosis include:
Prolonged exposure to naturally synthesized estrogen; for example, from late menopause or early menarche
Obstruction of menstrual outflow; for example, in Müllerian anomalies
Potential toxins:
Dioxins- Several studies have investigated the potential link between exposure to dioxins and endometriosis, but evidence is equivocal and potential mechanisms are poorly understood. A 2004 review of studies of dioxin and endometriosis concluded that "the human data supporting the dioxin-endometriosis association are scanty and conflicting", and a 2009 follow-up review also found that there was "insufficient evidence" in support of a link between dioxin exposure and developing endometriosis.
Endocrine-disrupting chemicals (EDCs)- A wider class of hormonally active agents, to which dioxin belongs, consists of both natural and manmade compounds, e.g., bisphenols, phthalates, pesticides (chlorpyrifos, hexachlorobenzene) and polychlorinated biphenyls (PCBs). Dietary uptake represents a significant source of EDC exposure via consumption of food, water and beverages, but exposure can also occur through ingestion of EDC dust and inhalation of its gases or particles in the air. Most EDCs are lipophilic, allowing them to bioaccumulate in adipose tissue (body fat) and increase in concentration. Bisphenol A (BPA), bisphenol S (BPS), phthalates, pesticides and PCBs all have a suspected linkage to endometriosis, though have not been definitively proven as being causative.
Vaginal dysbiosis
A growing body of evidence has shown a correlation between an imbalance in the vaginal microbiome and the appearance of endometriosis. This correlation is mediated by an immune system overload in the context of retrograde menstruation, in which it fails to detect and kill cells that come outside of the vaginal environment. By disrupting normal immune function, dysbiosis leads to elevated levels of proinflammatory cytokines, a compromised immunosurveillance system and altered immune cell profiles. Indeed, the activation of Toll-like receptors in macrophages leads to a greater activity of this immune cell type. They, in turn, secrete factors (such as the pro-inflammatory cytokine interleukin 8) that help creating an inflammatory environment, ultimately favoring the proliferation and adhesion of endometrial cells.
Pathophysiology
While the exact cause of endometriosis remains unknown, many theories have been presented to better understand and explain its development. These concepts do not necessarily exclude each other. The pathophysiology of endometriosis is likely to be multifactorial and to involve an interplay between several factors.
Formation
The main theories for the formation of the ectopic endometrium-like tissue include retrograde menstruation, Müllerianosis, coelomic metaplasia, vascular dissemination of stem cells, and surgical transplantation were postulated as early as 1870. Each is further described below.
Retrograde menstruation theory
The theory of retrograde menstruation (also called the implantation theory or transplantation theory) is the most commonly accepted theory for the dissemination and transformation of ectopic endometrium into endometriosis. It suggests that during a woman's menstrual flow, some of the endometrial debris flow backward through the fallopian tubes and into the peritoneal cavity, attaching itself to the peritoneal surface (the lining of the abdominal cavity) where it can proceed to invade the tissue as or transform into endometriosis. It is not clear at what stage the transformation of endometrium, or any cell of origin such as stem cells or coelomic cells (see those theories below), to endometriosis begins.
Proofs in support of the theory are based on retrospective epidemiological studies that an association with endometrial implants attached to the peritoneal cavity, which would develop into endometrial lesions and retrograde menstruation; and the fact that animals like rodents and non-human primates whose endometrium is not shed during the estrous cycle don't develop naturally endometriosis contrary to animals that have a natural menstrual cycle like rhesus monkeys and baboons.
Retrograde menstruation alone is not able to explain all instances of endometriosis, and additional factors such as genetics, immunology, stem cell migration, and coelomic metaplasia (see "Other theories" on this page) are needed to account for disseminated disease and why many individuals with retrograde menstruation are not diagnosed with endometriosis. In addition, endometriosis has shown up in people who have never experienced menstruation including cisgender men, fetuses, and prepubescent girls. Further theoretical additions are needed to complement the retrograde menstruation theory to explain why cases of endometriosis show up in the brain and lungs.
Researchers are investigating the possibility that the immune system may not be able to cope with the cyclic onslaught of retrograde menstrual fluid. In this context there is interest in studying the relationship of endometriosis to autoimmune disease, allergic reactions, and the impact of toxic materials. It is still unclear what, if any, causal relationship exists between toxic materials or autoimmune disease and endometriosis. There are immune system changes in people with endometriosis, such as an increase of macrophage-derived secretion products, but it is unknown if these are contributing to the disorder or are reactions from it.
Endometriotic lesions differ in their biochemistry, hormonal response, immunology, inflammatory response when compared to endometrium. This is likely because the cells that give rise to endometriosis are a side population of cells. Similarly, there are changes in, for example, the mesothelium of the peritoneum in people with endometriosis, such as loss of tight junctions, but it is unknown if these are causes or effects of the disorder.
In rare cases where imperforate hymen does not resolve itself prior to the first menstrual cycle and goes undetected, blood and endometrium are trapped within the uterus until such time as the problem is resolved by surgical incision. Many health care practitioners never encounter this defect, and due to the flu-like symptoms it is often misdiagnosed or overlooked until multiple menstrual cycles have passed. By the time a correct diagnosis has been made, endometrium and other fluids have filled the uterus and Fallopian tubes with results similar to retrograde menstruation resulting in endometriosis. The initial stage of endometriosis may vary based on the time elapsed between onset and surgical procedure.
The theory of retrograde menstruation as a cause of endometriosis was first proposed by John A. Sampson.
Other theories
Stem cells: Endometriosis may arise from stem cells from bone marrow and potentially other sources. In particular, this theory explains endometriosis found in areas remote from the pelvis such as the brain or lungs. Stem cells may be from local cells such as the peritoneum (see coelomic metaplasia below) or cells disseminated in the blood stream (see vascular dissemination below) such as those from the bone marrow.
Vascular dissemination: Vascular dissemination is a 1927 theory that has been revived with new studies of bone-marrow stem cells involved in pathogenesis.
Environment: Environmental toxins (e.g., dioxin, nickel) may cause endometriosis. Toxins such as dioxins and dioxin-like compounds tend to bioaccumulate within the human body. Further research is needed but "it is plausible that inflammatory-like processes, caused by dioxin-like environmental chemicals, can alter normal endometrial and immune cell physiology allowing persistence and development of endometrial tissue within the peritoneal cavity, normally cleared by immune system cells".
Müllerianosis: A theory supported by foetal autopsy is that cells with the potential to become endometrial, which are laid down in tracts during embryonic development called the female reproductive (Müllerian) tract as it migrates downward at 8–10 weeks of embryonic life, could become dislocated from the migrating uterus and act like seeds or stem cells.
Coelomic metaplasia: Coelomic cells which are the common ancestor of endometrial and peritoneal cells may undergo metaplasia (transformation) from one type of cell to the other, perhaps triggered by inflammation.
Vasculogenesis: Up to 37% of the microvascular endothelium of ectopic endometrial tissue originates from endothelial progenitor cells, which result in de novo formation of microvessels by the process of vasculogenesis rather than the conventional process of angiogenesis.
Neural growth: An increased expression of new nerve fibres is found in endometriosis but does not fully explain the formation of ectopic endometriotic tissue and is not definitely correlated with the amount of perceived pain.
Autoimmune: Graves disease is an autoimmune disease characterized by hyperthyroidism, goiter, ophthalmopathy, and dermopathy. People with endometriosis had higher rates of Graves disease. One of these potential links between Graves disease and endometriosis is autoimmunity.
Oxidative stress: Influx of iron is associated with the local destruction of the peritoneal mesothelium, leading to the adhesion of ectopic endometriotic cells. Peritoneal iron overload has been suggested to be caused by the destruction of erythrocytes, which contain the iron-binding protein hemoglobin, or a deficiency in the peritoneal iron metabolism system. Oxidative stress activity and reactive oxygen species (ROS) (such as superoxide anions and peroxide levels) are reported to be higher than normal in people with endometriosis. Oxidative stress and the presence of excess ROS can damage tissue and induce rapid cellular division. Mechanistically, there are several cellular pathways by which oxidative stress may lead to or may induce proliferation of endometriotic lesions, including the mitogen activated protein (MAP) kinase pathway and the extracellular signal-related kinase (ERK) pathway. Activation of both of the MAP and ERK pathways lead to increased levels of c-Fos and c-Jun, which are proto-oncogenes that are associated with high-grade lesions.
Localization
Most often, endometriosis is found on the:
Ovaries
Fallopian tubes
Tissues that hold the uterus in place (ligaments)
Outer surface of the uterus
Less common pelvic sites are:
Vagina
Cervix
Vulva
Bowel
Bladder
Rectum
Endometriosis may spread to the cervix and vagina or to sites of a surgical abdominal incision, known as "scar endometriosis." Rectovaginal or bowel endometriosis affects approximately 5-12% of those with endometriosis, and can cause severe pain with bowel movements.
Deep infiltrating endometriosis (DIE) has been defined as the presence of endometrial glands and stroma infiltrating more than 5 mm in the subperitoneal tissue. The prevalence of DIE is estimated to be 1 to 2% in women of reproductive age. Deep endometriosis typically presents as a single nodule in the vesicouterine fold or in the lower 20 cm of the bowel. Deep endometriosis can be associated with severe pain. However, it can be present without severe levels of pain.
Male endometriosis
Endometriosis has been reported in people assigned male at birth. Prostate endometriosis has been reported following estrogen therapy for prostate cancer and feminizing hormone therapy.
Abdominal endometriosis also happens in men following cirrhosis.
Extrapelvic endometriosis
Rarely, endometriosis appears in extrapelvic parts of the body, such as the lungs, brain, and skin. "Scar endometriosis" can occur in surgical abdominal incisions. Risk factors for scar endometriosis include previous abdominal surgeries, such as a hysterotomy or cesarean section, or ectopic pregnancies, salpingostomy puerperal sterilization, laparoscopy, amniocentesis, appendectomy, episiotomy, vaginal hysterectomies, and hernia repair.
Endometriosis may also present with skin lesions in cutaneous endometriosis.
Less commonly lesions can be found on the diaphragm or lungs. Diaphragmatic endometriosis is rare, almost always on the right hemidiaphragm, and may inflict the cyclic pain of the right scapula (shoulder) or cervical area (neck) during a menstrual period. Pulmonary endometriosis can be associated with a thoracic endometriosis syndrome that can include catamenial (occurs during menstruation) pneumothorax seen in 73% of women with the syndrome, catamenial hemothorax in 14%, catamenial hemoptysis in 7%, and pulmonary nodules in 6%.
Diagnosis
A health history and a physical examination can lead the health care practitioner to suspect endometriosis. There is a clear benefit for performing a transvaginal ultrasound (TVUS) as a first step of testing for endometriosis.
Definitive diagnosis is based on the morphology (form and structure) of the pelvic region, determined by observation (surgical or non-invasive imaging), classified into four different stages of endometriosis. The American Society of Reproductive Medicine's scale, revised in 1996, gives higher scores to deep, thick lesions or intrusions on the ovaries and dense, enveloping adhesions on the ovaries or fallopian tubes. Additionally, histological studies, when performed, should show specific findings.
For many patients, there are significant delays in diagnosis. Studies show an average delay of 11.7 years in the United States. Patients in the UK have an average delay of 8 years and in Norway of 6.7 years. A third of women had consulted their GP six or more times before being diagnosed.
The most common sites of endometriosis are the ovaries, followed by the Douglas pouch, the posterior leaves of the broad ligaments, and the sacrouterine ligaments.
As for deep infiltrating endometriosis, TVUS, TRUS and MRI are the techniques of choice for non-invasive diagnosis with a high sensitivity and specificity.
Laparoscopy
Laparoscopy, a surgical procedure where a camera is used to look inside the abdominal cavity, is the only way to accurately diagnose the extent and severity of pelvic/abdominal endometriosis. Laparoscopy is not an applicable test for extrapelvic sites such as umbilicus, hernia sacs, abdominal wall, lung, or kidneys.
Reviews in 2019 and 2020 concluded that 1) with advances in imaging, endometriosis diagnosis should no longer be considered synonymous with immediate laparoscopy for diagnosis, and 2) endometriosis should be classified a syndrome that requires confirmation of visible lesions seen at laparoscopy in addition to characteristic symptoms.
Laparoscopy permits lesion visualization unless the lesion is visible externally (e.g., an endometriotic nodule in the vagina) or is extra-abdominal. If the growths (lesions) are not visible, a biopsy must be taken to determine the diagnosis. Surgery for diagnoses also allows for surgical treatment of endometriosis at the same time.
During a laparoscopic procedure, lesions can appear dark blue, powder-burn black, red, white, yellow, brown or non-pigmented. Lesions vary in size. Some within the pelvis walls may not be visible, as normal-appearing peritoneum of infertile women reveals endometriosis on biopsy in 6–13% of cases. Early endometriosis typically occurs on the surfaces of organs in the pelvic and intra-abdominal areas. Health care providers may call areas of endometriosis by different names, such as implants, lesions, or nodules. Larger lesions may be seen within the ovaries as endometriomas or "chocolate cysts", "chocolate" because they contain a thick brownish fluid, mostly old blood.
Frequently during diagnostic laparoscopy, no lesions are found in individuals with chronic pelvic pain, a symptom common to other disorders including adenomyosis, pelvic adhesions, pelvic inflammatory disease, congenital anomalies of the reproductive tract, and ovarian or tubal masses.
Ultrasound
Vaginal ultrasound can be used to diagnosis endometriosis, or for localizing endometrioma before surgery. This can be used to identify the spread of disease in individuals with well-established clinical suspicion of endometriosis. Vaginal ultrasound is inexpensive, easily accessible, has no contraindications and requires no preparation. By extending the ultrasound assessment into the posterior and anterior pelvic compartments a sonographer is able to evaluate structural mobility and look for deep infiltrating endometriotic nodules. Better sonographic detection of deep infiltrating endometriosis could reduce the number of diagnostic laparoscopies, as well as guide disease management and enhance patient quality of life.
Magnetic resonance imaging
MRI is another means of detecting lesions in a non-invasive manner. MRI is not widely used due to its cost and limited availability, although it can be used to detect the most common form of endometriosis (endometrioma) with a sufficient accuracy. A 2020 article recommended administering an anti-spasmodic agent (i.e. hyoscine butylbromide) and a big glass of water (if the bladder is empty), and scanning in the supine position with an abdominal strap, for better image quality. It also recommended using pelvic-phased array coils and T1 (spin-lattice) weighted scanning, with and without suppression of fat for endometriomas, and sagittal, axial and oblique 2D T2 (spin-spin) weighting for deep infiltrating endometriosis.
Stages of disease
By surgical observation, endometriosis can be classified as stage I–IV by the 1996 scale of the American Society of Reproductive Medicine (ASRM). The scale uses a point system that assesses lesions and adhesions in the pelvic organs. It is important to note that staging assesses physical disease only, not the level of pain or infertility. A person with Stage I endometriosis may have a little disease and severe pain, while a person with Stage IV endometriosis may have severe disease and no pain or vice versa. The various stages are summarized by:
Stage I (Minimal)
Findings restricted to only superficial lesions and possibly a few filmy adhesions.
Stage II (Mild)
In addition, some deep lesions are present in the cul-de-sac.
Stage III (Moderate)
As above, plus the presence of endometriomas on the ovary and more adhesions.
Stage IV (Severe)
As above, plus large endometriomas, extensive adhesions. Implants and adhesions may be found beyond the uterus. Large ovarian cysts are common.
Markers
An area of research is the search for endometriosis markers.
In 2010, essentially all proposed biomarkers for endometriosis were of unclear medical use, although some appear to be promising. The one biomarker that has been in use over the last 20 years is CA-125. A 2016 review found that this biomarker was present in those with symptoms of endometriosis; and, once ovarian cancer has been ruled out, a positive CA-125 may confirm the diagnosis. Its performance in ruling out endometriosis is low. CA-125 levels appear to fall during endometriosis treatment, but it has not shown a correlation with disease response.
Another review in 2011 identified several putative biomarkers upon biopsy, including findings of small sensory nerve fibers or defectively expressed β3 integrin subunit. It has been postulated a future diagnostic tool for endometriosis will consist of a panel of several specific and sensitive biomarkers, including both substance concentrations and genetic predisposition.
A 2016 review of endometrial biomarkers for diagnosing endometriosis was unable to draw conclusions due to the low quality of the evidence.
MicroRNAs have the potential to be used in diagnostic and therapeutic decisions.
Histopathology
For a histopathological diagnosis, at least two of the following three criteria should be present:
Endometrial type stroma
Endometrial epithelium with glands
Evidence of chronic hemorrhage, mainly hemosiderin deposits
Immunohistochemistry has been found to be useful in diagnosing endometriosis as stromal cells have a peculiar surface antigen, CD10, thus allowing the pathologist go straight to a staining area and confirm the presence of stromal cells and sometimes glandular tissue is identified that was missed on routine H&E staining.
Pain quantification
The most common pain scale for quantification of endometriosis-related pain is the visual analogue scale (VAS); VAS and numerical rating scale (NRS) were the best adapted pain scales for pain measurement in endometriosis. For research purposes, and for more detailed pain measurement in clinical practice, VAS or NRS for each type of typical pain related to endometriosis (dysmenorrhea, deep dyspareunia and non-menstrual chronic pelvic pain), combined with the clinical global impression (CGI) and a quality of life scale, are used.
Prevention
Limited evidence indicates that the use of combined oral contraceptives is associated with a reduced risk of endometriosis, as is regular exercise and the avoidance of alcohol and caffeine. There is little known information on preventing endometriosis.
Management
While there is no cure for endometriosis, there are treatments for pain and endometriosis-associated infertility. Pain can be treated with hormones, painkillers or, in severe cases, surgery.
In most cases, the symptoms disappear or improve with menopause (natural or surgical). In the reproductive years, endometriosis is merely managed: the goal is to provide pain relief, to restrict progression of the process, and to restore or preserve fertility where needed. In younger individuals, some surgical treatment attempts to remove endometriotic tissue and preserve the ovaries without damaging normal tissue.
Pharmacotherapy for pain management can be initiated based on the presence of symptoms and examination and ultrasound findings that rule out other potential causes.
In general, the diagnosis of endometriosis is confirmed during surgery, at which time removal can be performed. Further steps depend on circumstances: someone without infertility can manage symptoms with pain medication and hormonal medication that suppresses the natural cycle, while an infertile individual may be treated expectantly after surgery, with fertility medication, or with in vitro fertilisation (IVF).
A 2020 Cochrane systematic review found that for all types of endometriosis, "it is uncertain whether laparoscopic surgery improves overall pain compared to diagnostic laparoscopy".
Surgery
Based on strong evidence, experts recommend that surgery be performed laparoscopically (through keyhole surgery) rather than open. Treatment consists of the ablation or excision of the endometriosis, electrocoagulation, lysis of adhesions, resection of endometriomas, and restoration of normal pelvic anatomy as much as is possible. When laparoscopic surgery is used, small instruments are inserted through the incisions to remove the endometriosis tissue and adhesions. Because the incisions are very small, there will only be small scars on the skin after the procedure, and most individuals recover from surgery quickly and have a reduced risk of adhesions. Many endometriosis specialists believe that excision is the ideal surgical method to treat endometriosis. A 2017 literature review found excision improved some outcomes over ablation. In the United States, some specialists trained in excision for endometriosis do not accept health insurance, because insurance companies do not reimburse the higher costs of this procedure over ablation.
As for deep endometriosis, a segmental resection or shaving of nodules is effective but is associated with an increased rate of complications, of which about 4.6% are major.
Historically, a hysterectomy (removal of the uterus) was thought to be a cure for endometriosis in individuals who do not wish to conceive. Removal of the uterus may be beneficial as part of the treatment, if the uterus itself is affected by adenomyosis. However, this should only be done in combination with removal of the endometriosis by excision. If endometriosis is not also removed at the time of hysterectomy, pain may persist. A study of hysterectomy patients found those with endometriosis were not using less pain medication 3 years after the procedure.
Presacral neurectomy may be performed where the nerves to the uterus are cut. However, this technique is not usually used due to the high incidence of associated complications including presacral hematoma and irreversible problems with urination and constipation.
Recurrence
The underlying process that causes endometriosis may not cease after a surgical or medical intervention. Even though surgery can improve symptoms, the resurgence of pain is common. A study has shown that dysmenorrhea recurs at a rate of 30 percent within a year following laparoscopic surgery. Resurgence of lesions tend to appear in the same location if the lesions were not completely removed during surgery. It has been shown that laser ablation resulted in higher and earlier recurrence rates when compared with endometrioma cystectomy; and recurrence after repetitive laparoscopy was similar to that after the first surgery. Endometriosis has a 10% recurrence rate after hysterectomy and bilateral salpingo-oophorectomy.
Endometriosis recurrence following conservative surgery is estimated as 21.5% at 2 years and 40-50% at 5 years.
Recurrence rate for DIE after surgery is less than 1%.
Risks and safety of pelvic surgery
Risk of developing complications following surgery depend on the type of the lesion that has undergone surgery. 55% to 100% of individuals develop adhesions following pelvic surgery, which can result in infertility, chronic abdominal and pelvic pain, and difficult reoperative surgery. Trehan's temporary ovarian suspension, a technique in which the ovaries are suspended for a week after surgery, may be used to reduce the incidence of adhesions after endometriosis surgery. Removal of cysts on the ovary without removing the ovary is a safe procedure.
Hormonal medications
Hormonal birth control therapy: Birth control pills reduce the menstrual pain and recurrence rate for endometrioma following conservative surgery for endometriosis. A 2018 Cochrane systematic review found that there is insufficient evidence to make a judgement on the effectiveness of the combined oral contraceptive pill compared with placebo or other medical treatment for managing pain associated with endometriosis partly because of lack of included studies for data analysis (only two for COCP vs placebo).
Progestogens: Progesterone counteracts estrogen and inhibits the growth of the endometrium. Danazol and gestrinone are suppressive steroids with some androgenic activity. Both agents inhibit the growth of endometriosis but their use has declined, due in part to virilizing side effects such as excessive hair growth and voice changes. There is tentative evidence based on cohort studies that dienogest and norethisterone acetate (NETA) may help patients with DIE in terms of pain. There is tentative evidence based on a prospective study that vaginal danazol reduces pain in those affected by DIE.
Gonadotropin-releasing hormone (GnRH) modulators: These drugs include GnRH agonists such as leuprorelin and GnRH antagonists such as elagolix and are thought to work by decreasing estrogen levels. A 2010 Cochrane review found that GnRH modulators were more effective for pain relief in endometriosis than no treatment or placebo, but were not more effective than danazol or intrauterine progestogen, and had more side effects than danazol. A 2018 Swedish systematic review found that GnRH modulators had similar pain-relieving effects to gestagen, but also decreased bone density.
Aromatase inhibitors are medications that block the formation of estrogen and have become of interest for researchers who are treating endometriosis. Examples of aromatase inhibitors include anastrozole and letrozole. Evidence for aromatase inhibitors is confirmed by numerous controlled studies that show benefit in terms of pain control and quality of life when used in combination with gestagens or oral contraceptives with less side-effects when used in combination with oral contraceptives like norethisterone acetate. Despite multiple benefits, there are lot of things to consider before using aromatase inhibitors for endometriosis, as it is common for them to induce functional cysts as an adverse effects. Moreover, dosages, treatment length, appropriate add-back therapies and mode of administration is still being investigated.
Progesterone receptor modulators like mifepristone and gestrinone have the potential (based on only one randomized controlled trial each) to be used as a treatment to manage pain caused by endometriosis.
Other medicines
Melatonin, there is tentative evidence for its use (at a dose of 10 mg) in reducing pain related to endometriosis.
Vitamin C and E antioxidant supplementation may significantly reduce pain symptoms in endometriosis.
Opioids: Morphine sulphate tablets and other opioid painkillers work by mimicking the action of naturally occurring pain-reducing chemicals called "endorphins". There are different long acting and short acting medications that can be used alone or in combination to provide appropriate pain control.
Chinese herbal medicine was reported to have comparable benefits to gestrinone and danazol in patients who had had laparoscopic surgery, though the review notes that the two trials were small and of "poor methodological quality" and results should be "interpreted cautiously" as better quality research is needed.
Serrapeptase, a digestive enzyme found in the intestines of silkworms. Serrapeptase is widely used in Japan and Europe as an anti-inflammatory treatment. More research is needed but serrapeptase may be used by endometriosis patients to reduce inflammation.
Angiogenesis inhibitors lack clinical evidence of efficacy in endometriosis therapy. Under experimental in vitro and in vivo conditions, compounds that have been shown to exert inhibitory effects on endometriotic lesions include growth factor inhibitors, endogenous angiogenesis inhibitors, fumagillin analogues, statins, cyclo-oxygenase-2 inhibitors, phytochemical compounds, immunomodulators, dopamine agonists, peroxisome proliferator-activated receptor agonists, progestins, danazol and gonadotropin-releasing hormone agonists. However, many of these agents are associated with undesirable side effects and more research is necessary. An ideal therapy would diminish inflammation and underlying symptoms without being contraceptive.
Pentoxifylline, an immunomodulating agent, has been theorized to improve pain as well as improve pregnancy rates in individuals with endometriosis. There is not enough evidence to support the effectiveness or safety of either of these uses. Current American Congress of Obstetricians and Gynecologists (ACOG) guidelines do not include immunomodulators, such as pentoxifylline, in standard treatment protocols.
NSAIDs are anti-inflammatory medications commonly used for endometriosis patients despite unproven efficacy and unintended adverse effects.
Neuromodulators like gabapentin did not prove to be superior to placebo in managing pain caused by endometriosis.
The overall effectiveness of manual physical therapy to treat endometriosis has not yet been identified.
Comparison of interventions
A 2021 meta-analysis found that GnRH analogues and combined hormonal contraceptives were the best treatment for reducing dyspareunia, menstrual and non menstrual pelvic pain. A 2018 Swedish systematic review found a large number of studies but a general lack of scientific evidence for most treatments. There was only one study of sufficient quality and relevance comparing the effect of surgery and non-surgery. Cohort studies indicate that surgery is effective in decreasing pain. Most complications occurred in cases of low intestinal anastomosis, while risk of fistula occurred in cases of combined abdominal or vaginal surgery, and urinary tract problems were common in intestinal surgery. The evidence was found to be insufficient regarding surgical intervention.
The advantages of physical therapy techniques are decreased cost, absence of major side-effects, it does not interfere with fertility, and near-universal increase of sexual function. Disadvantages are that there are no large or long-term studies of its use for treating pain or infertility related to endometriosis.
Treatment of infertility
Surgery is more effective than medicinal intervention for addressing infertility associated with endometriosis. Surgery attempts to remove endometrium-like tissue and preserve the ovaries without damaging normal tissue. Receiving hormonal suppression therapy after surgery might be positive regarding endometriosis recurrence and pregnancy. In-vitro fertilization (IVF) procedures are effective in improving fertility in many individuals with endometriosis.
During fertility treatment, the ultralong pretreatment with GnRH-agonist has a higher chance of resulting in pregnancy for individuals with endometriosis, compared to the short pretreatment.
Epidemiology
Determining how many people have endometriosis is challenging because definitive diagnosis requires surgical visualization through laparoscopic surgery. Criteria that are commonly used to establish a diagnosis include pelvic pain, infertility, surgical assessment, and in some cases, magnetic resonance imaging. An ultrasound can identify large clumps of tissue as potential endometriosis lesions and ovarian cysts but it is not effective for all patients, especially in cases with smaller, superficial lesions.
Ethnic differences in endometriosis have been observed. The condition is more common in women of East Asian and Southeast Asian descent than in White women. Risk factors include having a family history of the condition.
One estimate is that 10.8 million people are affected globally . Other sources estimate 6 to 10% of the general female population and 2 to 11% of asymptomatic women are affected. In addition, 11% of women in a general population have undiagnosed endometriosis that can be seen on magnetic resonance imaging (MRI). Endometriosis is most common in those in their thirties and forties; however, it can begin in girls as early as eight years old. It results in few deaths with unadjusted and age-standardized death rates of 0.1 and 0.0 per 100,000. Endometriosis was first determined to be a separate condition in the 1920s. Before that time, endometriosis and adenomyosis were considered together. It is unclear who first described the disease.
It chiefly affects adults from premenarche to postmenopause, regardless of race or ethnicity or whether or not they have had children and is estimated to affect over 190 million women in their reproductive years. Incidences of endometriosis have occurred in postmenopausal individuals, and in less common cases, individuals may have had endometriosis symptoms before they even reach menarche.
The rate of recurrence of endometriosis is estimated to be 40-50% for adults over a 5-year period. The rate of recurrence has been shown to increase with time from surgery and is not associated with the stage of the disease, initial site, surgical method used, or post-surgical treatment.
History
Endometriosis was first discovered microscopically by Karl von Rokitansky in 1860, although the earliest antecedents may have stemmed from concepts published almost 4,000 years ago. The Hippocratic Corpus outlines symptoms similar to endometriosis, including uterine ulcers, adhesions, and infertility. Historically, women with these symptoms were treated with leeches, straitjackets, bloodletting, chemical douches, genital mutilation, pregnancy (as a form of treatment), hanging upside down, surgical intervention, and even killing due to suspicion of demonic possession. Hippocratic doctors recognized and treated chronic pelvic pain as a true organic disorder 2,500 years ago, but during the Middle Ages, there was a shift into believing that women with pelvic pain were mad, immoral, imagining the pain, or simply misbehaving. The symptoms of inexplicable chronic pelvic pain were often attributed to imagined madness, female weakness, promiscuity, or hysteria. The historical diagnosis of hysteria, which was thought to be a psychological disease, may have indeed been endometriosis. The idea that chronic pelvic pain was related to mental illness influenced modern attitudes regarding individuals with endometriosis, leading to delays in correct diagnosis and indifference to the patients' true pain throughout the 20th and into the 21st century.
Hippocratic doctors believed that delaying childbearing could trigger diseases of the uterus, which caused endometriosis-like symptoms. Women with dysmenorrhea were encouraged to marry and have children at a young age. The fact that Hippocratics were recommending changes in marriage practices due to an endometriosis-like illness implies that this disease was likely common, with rates higher than the 5-15% prevalence that is often cited today. If indeed this disorder was so common historically, this may point away from modern theories that suggest links between endometriosis and dioxins, PCBs, and chemicals.
The early treatment of endometriosis was surgical and included oophorectomy (removal of the ovaries) and hysterectomy (removal of the uterus). In the 1940s, the only available hormonal therapies for endometriosis were high-dose testosterone and high-dose estrogen therapy. High-dose estrogen therapy with diethylstilbestrol for endometriosis was first reported by Karnaky in 1948 and was the main pharmacological treatment for the condition in the early 1950s. Pseudopregnancy (high-dose estrogen–progestogen therapy) for endometriosis was first described by Kistner in the late 1950s. Pseudopregnancy as well as progestogen monotherapy dominated the treatment of endometriosis in the 1960s and 1970s. These agents, although efficacious, were associated with intolerable side effects. Danazol was first described for endometriosis in 1971 and became the main therapy in the 1970s and 1980s. In the 1980s GnRH agonists gained prominence for the treatment of endometriosis and by the 1990s had become the most widely used therapy. Oral GnRH antagonists such as elagolix were introduced for the treatment of endometriosis in 2018.
Society and culture
Public figures
A number of public figures have spoken about their experience with endometriosis, including:
RuthAnne
Emma Barnett
Emma Bunton
Alexa Chung
Danielle Collins
Olivia Culpo
Lena Dunham
Abby Finkenauer
Bethenny Frankel
Whoopi Goldberg
Mel Greig
Halsey
Emma Hayes
Julianne Hough
Bridget Hustwaite
Bindi Irwin
Jaime King
Padma Lakshmi
Cyndi Lauper
Jillian Michaels
Monica
Marilyn Monroe
Tia Mowry
Sinéad O'Connor
Dolly Parton
Florence Pugh
Daisy Ridley
Emma Roberts
Susan Sarandon
Amy Schumer
Kirsten Storms
Gabrielle Union
Lacey Schwimmer
Chrissy Teigen
Emma Watkins
Mae Whitman
Jessica Williams
Leah Williamson
Economic burden
The economic burden of endometriosis is widespread and multifaceted. Endometriosis is a chronic disease that has direct and indirect costs which include loss of work days, direct costs of treatment, symptom management, and treatment of other associated conditions such as depression or chronic pain. One factor which seems to be associated with especially high costs is the delay between onset of symptoms and diagnosis.
Costs vary greatly between countries. Two factors that contribute to the economic burden include healthcare costs and losses in productivity. A Swedish study of 400 endometriosis patients found "Absence from work was reported by 32% of the women, while 36% reported reduced time at work because of endometriosis". An additional cross sectional study with Puerto Rican women, "found that endometriosis-related and coexisting symptoms disrupted all aspects of women's daily lives, including physical limitations that affected doing household chores and paid employment. The majority of women (85%) experienced a decrease in the quality of their work; 20% reported being unable to work because of pain, and over two-thirds of the sample continued to work despite their pain."
Medical culture
There are a number of barriers that those affected face to receiving diagnosis and treatment for endometriosis. Some of these include outdated standards for laparoscopic evaluation, stigma about discussing menstruation and sex, lack of understanding of the disease, primary-care physicians' lack of knowledge, and assumptions about typical menstrual pain. On average, those later diagnosed with endometriosis waited 2.3 years after the onset of symptoms before seeking treatment and nearly three quarters of women receive a misdiagnosis prior to endometriosis. Self-help groups say practitioners delay making the diagnosis, often because they do not consider it a possibility. There is a typical delay of 7–12 years from symptom onset in affected individuals to professional diagnosis. There is a general lack of knowledge about endometriosis among primary care physicians. Half of general health care providers surveyed in a 2013 study were unable to name three symptoms of endometriosis. Health care providers are also likely to dismiss described symptoms as normal menstruation. Younger patients may also feel uncomfortable discussing symptoms with a physician.
Race and ethnicity
Race and ethnicity may play a role in how endometriosis affects one's life. Endometriosis is less thoroughly studied among Black people, and the research that has been done is outdated.
Cultural differences among ethnic groups also contribute to attitudes toward and treatment of endometriosis, especially in Hispanic or Latino communities. A study done in Puerto Rico in 2020 found that health care and interactions with friends and family related to discussing endometriosis were affected by stigma. The most common finding was a referral to those expressing pain related to endometriosis as "changuería" or "changas", terms used in Puerto Rico to describe pointless whining and complaining, often directed at children.
Stigma
The existing stigma surrounding women's health, specifically endometriosis, can lead to patients not seeking diagnoses, lower quality of healthcare, increased barriers to care and treatment, and negative reception from members of society. Additionally, menstrual stigma significantly contributes to the broader issue of endometriosis stigma, creating an interconnected challenge that extends beyond reproductive health.
Widespread awareness campaigns, developments and implementations aimed to multilevel anti-stigma organizational and structural changes, as well as more qualitative studies of the endometriosis stigma, help to overcome the harm of the phenomenon.
Research
Preliminary research on mouse models showed that monoclonal antibodies, as well as inhibitors of MyD88 downstream signaling pathway, can reduce lesion volume. Thanks to that, clinical trials are being done on using a monoclonal antibody directed against IL-33 and using anakinra, an IL-1 receptor antagonist.
Taking the contraceptive pills or getting long-acting progestogen injections seems to be equally effective for preventing recurring pain after endometriosis surgery. Compared to taking the pill, progestogen might result in reduced risk of needing further treatments or surgery.
Clinical trials are exploring the potential benefits of cannabinoid extracts, dichloroacetic acid and curcuma capsules.
| Biology and health sciences | Specific diseases | Health |
53338 | https://en.wikipedia.org/wiki/Cervical%20cancer | Cervical cancer | Cervical cancer is a cancer arising from the cervix or in any layer of the wall of the cervix. It is due to the abnormal growth of cells that have the ability to invade or spread to other parts of the body. Early on, typically no symptoms are seen. Later symptoms may include abnormal vaginal bleeding, pelvic pain or pain during sexual intercourse. While bleeding after sex may not be serious, it may also indicate the presence of cervical cancer.
Virtually all cervical cancer cases (99%) are linked to genital human papillomavirus infection (HPV); most who have had HPV infections, however, do not develop cervical cancer. HPV 16 and 18 strains are responsible for approximately 70% of cervical cancer cases globally and nearly 50% of high grade cervical pre-cancers. Minor risk factors include smoking, a weak immune system, birth control pills, starting sex at a young age, and having many sexual partners. Genetic factors also contribute to cervical cancer risk. Cervical cancer typically develops from precancerous changes called cervical intraepithelial neoplasia over 10 to 20 years. About 90% of cervical cancer cases are squamous cell carcinomas, 10% are adenocarcinoma, and a small number are other types. Diagnosis is typically by cervical screening followed by a biopsy. Medical imaging is then done to determine whether or not the cancer has spread.
HPV vaccination is the most cost-effective public health measure against cervical cancer. There are six licensed HPV vaccines. They protect against two to seven high-risk strains of this family of viruses and may prevent up to 90% of cervical cancers. By the end of 2023, 143 countries (74% of WHO member states) provided the HPV vaccine in their national immunization schedule for girls. As of 2022, 47 countries (24% of WHO member states) also did it for boys. As a risk of cancer still exists, guidelines recommend continuing regular Pap tests. Other methods of prevention include having few or no sexual partners and the use of condoms. Cervical cancer screening using the Pap test or acetic acid can identify precancerous changes, which when treated, can prevent the development of cancer. Treatment may consist of some combination of surgery, chemotherapy, and radiation therapy. Five-year survival rates in the United States are 68%. Outcomes, however, depend very much on how early the cancer is detected.
Worldwide, cervical cancer is both the fourth-most common type of cancer and the fourth-most common cause of death from cancer in women, with over 660,000 new cases and around 350,000 deaths in 2022. This is about 8% of the total cases and total deaths from cancer. 88% (2020 figure) of cervical cancers and 90% of deaths occur in low- and middle-income countries and 2% (2020 figure) in high-income countries. Of the 20 hardest hit countries by cervical cancer, 19 are in Africa. In low-income countries, it is one of the most common causes of cancer death with an incidence rate of 47.3 per 100,000 women. In developed countries, the widespread use of cervical screening programs has dramatically reduced rates of cervical cancer. Expected scenarios for the reduction of mortality due to cervical cancer worldwide (and specially in low-income countries) have been reviewed, given assumptions with respect to the achievement of recommended prevention targets using triple-intervention strategies defined by WHO. In medical research, the most famous immortalized cell line, known as HeLa, was developed from cervical cancer cells of a woman named Henrietta Lacks.
17 November is the Cervical Cancer Elimination Day of Action. The date marks the day in 2020 when WHO launched the Global strategy to accelerate the elimination of cervical cancer as a public health problem, with a resolution passed by 194 countries. To eliminate cervical cancer, all countries must reach and maintain an incidence rate of below 4 per 100 000 women.
Signs and symptoms
The early stages of cervical cancer may be completely free of symptoms. Vaginal bleeding, contact bleeding (one most common form being bleeding after sexual intercourse), or (rarely) a vaginal mass may indicate the presence of malignancy. Also, moderate pain during sexual intercourse and vaginal discharge are symptoms of cervical cancer. In advanced disease, metastases may be present in the abdomen, lungs, or elsewhere.
Symptoms of advanced cervical cancer may include loss of appetite, weight loss, fatigue, pelvic pain, back pain, leg pain, swollen legs, heavy vaginal bleeding, bone fractures, and (rarely) leakage of urine or faeces from the vagina. Bleeding after douching or after a pelvic exam is a common symptom of cervical cancer.
Causes
Infection with some types of HPV is the greatest risk factor for cervical cancer, followed by smoking. HIV infection is also a risk factor. Not all of the causes of cervical cancer are known, however, and several other contributing factors have been implicated.
Human papillomavirus
Infection with HPV is generally believed to be required for cervical cancer to occur. HPV types 16 and 18 are the cause of 75% of cervical cancer cases globally, while 31 and 45 are the causes of another 10%.
Women who have multiple sexual partners, or have partners who have multiple sexual partners, regardless of sex are at higher risk of cervical cancer.
Over 200 types of HPV known, 12 are classified as high-risk types (16, 18, 31, 33, 35, 39, 45, 51, 52, 56, 58, and 59), three as probable high-risk (26, 53, and 66), and 12 as low-risk (6, 11, 40, 42, 43, 44, 54, 61, 70, 72, 81, and CP6108).
Genital warts, which are a form of benign tumor of epithelial cells, are also caused by various strains of HPV. However, these serotypes are usually not related to cervical cancer. Having multiple strains at the same time is common, including those that can cause cervical cancer along with those that cause warts.
Smoking
Cigarette smoking, both active and passive, increases the risk of cervical cancer. Among HPV-infected women, current and former smokers have roughly two to three times the incidence of invasive cancer. Passive smoking is also associated with increased risk, but to a lesser extent.
Smoking has also been linked to the development of cervical cancer. Smoking can increase the risk in women a few different ways, which can be by direct and indirect methods of inducing cervical cancer. A direct way of contracting this cancer is a smoker has a higher chance of cervical intraepithelial neoplasia (CIN3) occurring, which has the potential of forming cervical cancer. When CIN3 lesions lead to cancer, most of them have the assistance of the HPV virus, but that is not always the case, which is why it can be considered a direct link to cervical cancer. Heavy smoking and long-term smoking seem to have more of a risk of getting the CIN3 lesions than lighter smoking or not smoking at all. Although smoking has been linked to cervical cancer, it aids in the development of HPV, which is the leading cause of this type of cancer. Also, not only does it aid in the development of HPV, but also if the woman is already HPV-positive, she is at an even greater likelihood of contracting cervical cancer.
Oral contraceptives
Long-term use of oral contraceptives is associated with increased risk of cervical cancer in women who have had HPV. Women who have used oral contraceptives for 5 to 9 years have about three times the incidence of invasive cancer, and those who used them for 10 years or longer have about four times the risk.
Multiple pregnancies
Having many pregnancies is associated with an increased risk of cervical cancer. Among HPV-infected women, those who have had seven or more full-term pregnancies have around four times the risk of cancer compared with women with no pregnancies, and two to three times the risk of women who have had one or two full-term pregnancies.
Diagnosis
Biopsy
The Pap test can be used as a screening test, but produces a false negative in up to 50% of cases of cervical cancer. Other concerns is the cost of doing Pap tests, which make them unaffordable in many areas of the world.
Confirmation of the diagnosis of cervical cancer or precancer requires a biopsy of the cervix. This is often done through colposcopy, a magnified visual inspection of the cervix aided by using a dilute acetic acid (e.g. vinegar) solution to highlight abnormal cells on the surface of the cervix, with visual contrast provided by staining the normal tissues a mahogany brown with Lugol's iodine. Medical devices used for biopsy of the cervix include punch forceps. Colposcopic impression, the estimate of disease severity based on the visual inspection, forms part of the diagnosis. Further diagnostic and treatment procedures are loop electrical excision procedure and cervical conization, in which the inner lining of the cervix is removed to be examined pathologically. These are carried out if the biopsy confirms severe cervical intraepithelial neoplasia.
Often before the biopsy, the doctor asks for medical imaging to rule out other causes of woman's symptoms. Imaging modalities such as ultrasound, CT scan, and MRI have been used to look for alternating disease, spread of the tumor, and effect on adjacent structures. Typically, they appear as heterogeneous mass on the cervix.
Interventions such as playing music during the procedure and viewing the procedure on a monitor can reduce the anxiety associated with the examination.
Precancerous lesions
Cervical intraepithelial neoplasia (CIN) means the development of abnormal cells on the surface of the cervix. It is caused by an HPV infection but in most cases it is resolved by the immune system. However a small percentage of people might develop a more serious CIN which, if left untreated, can develop into cervical cancer. CIN is often diagnosed during routine Pap smear examination or colposcopy.
The naming and histologic classification of cervical carcinoma precursor lesions has changed many times over the 20th century. The World Health Organization classification system was descriptive of the lesions, naming them mild, moderate, or severe dysplasia or carcinoma in situ (CIS). The term cervical intraepithelial neoplasia (CIN) was developed to place emphasis on the spectrum of abnormality in these lesions, and to help standardize treatment. For premalignant dysplastic changes, cervical intraepithelial neoplasia grading (CIN 1–3) is used. It classifies mild dysplasia as CIN1, moderate dysplasia as CIN2, and severe dysplasia and CIS as CIN3. More recently, CIN2 and CIN3 have been combined into CIN2/3. These results are what a pathologist might report from a biopsy.
These should not be confused with the Bethesda system terms for Pap test (cytopathology) results. Among the Bethesda results: Low-grade squamous intraepithelial lesion (LSIL) and high-grade squamous intraepithelial lesion (HSIL). An LSIL Pap may correspond to CIN1, and HSIL may correspond to CIN2 and CIN3, but they are results of different tests, and the Pap test results need not match the histologic findings.
Cancer subtypes
Histologic subtypes of invasive cervical carcinoma include:
Squamous cell carcinoma (about 80–85%)
adenocarcinoma (about 15% of cervical cancers in the UK)
Adenosquamous carcinoma
Small cell carcinoma
Neuroendocrine tumour
Glassy cell carcinoma
Villoglandular adenocarcinoma
Though squamous cell carcinoma is the cervical cancer with the most incidence, the incidence of adenocarcinoma of the cervix has been increasing in recent decades. Endocervical adenocarcinoma represents 20–25% of the histological types of cervical carcinoma. Gastric-type mucinous adenocarcinoma of the cervix is a rare type of cancer with aggressive behavior. This type of malignancy is not related to high-risk human papillomavirus (HPV).
Noncarcinoma malignancies which can rarely occur in the cervix include melanoma and lymphoma. The International Federation of Gynecology and Obstetrics (FIGO) stage does not incorporate lymph node involvement in contrast to the TNM staging for most other cancers. For cases treated surgically, information obtained from the pathologist can be used in assigning a separate pathologic stage, but is not to replace the original clinical stage.
Staging
Cervical cancer is staged by the FIGO system, which is based on clinical examination rather than surgical findings. Prior to the 2018 revisions to FIGO staging, the system allowed only these diagnostic tests to be used in determining the stage: palpation, inspection, colposcopy, endocervical curettage, hysteroscopy, cystoscopy, proctoscopy, intravenous urography, and X-ray examination of the lungs and skeleton, and cervical conization. However, the system allows use of any imaging or pathological methods for staging.
Prevention
Screening
Checking cervical cells with the Papanicolaou test (Pap test) for cervical pre-cancer has dramatically reduced the number of cases of, and mortality from, cervical cancer. Liquid-based cytology may reduce the number of inadequate samples. Pap test screening every three to five years with appropriate follow-up can reduce cervical cancer incidence up to 80%.
Pap test screening can reveal abnormal cells on the surface of the cervix called cervical intraepithelial neoplasia (CIN) that in a small percentage can develop into cervical cancer. These precancerous changes can be confirmed with further examination known as colposcopy.
Personal invitations encouraging women to get screened are effective at increasing the likelihood they will do so. Educational materials also help increase the likelihood women will go for screening, but they are not as effective as invitations.
According to the 2010 European guidelines, the age at which to start screening ranges between 20 and 30 years of age, but preferentially not before age 25 or 30 years, and depends on burden of the disease in the population and the available resources.
In the United States, screening is recommended to begin at age 21, regardless of age at which a woman began having sex or other risk factors. Pap tests should be done every three years between the ages of 21 and 65. In women over the age of 65, screening may be discontinued if no abnormal screening results were seen within the previous 10 years and no history of CIN2 or higher exists. HPV vaccination status does not change screening rates.
A number of recommended options exist for screening those 30 to 65. This includes cervical cytology every 3 years, HPV testing every 5 years, or HPV testing together with cytology every 5 years. Screening is not beneficial before age 25, as the rate of disease is low. Screening is not beneficial in women older than 60 years if they have a history of negative results. The American Society of Clinical Oncology guideline has recommend for different levels of resource availability.
Pap tests have not been as effective in developing countries. This is in part because many of these countries have an impoverished health care infrastructure, too few trained and skilled professionals to obtain and interpret Pap tests, uninformed women who get lost to follow-up, and a lengthy turn-around time to get results. Visual inspection with acetic acid and HPV DNA testing have been tried, though with mixed success.
Vaccination
There are six licensed HPV vaccines: three bivalent (protect against two types of HPV), two quadrivalent (against four), and one nonavalent vaccine (against nine) Three HPV vaccines (Gardasil, Gardasil 9, and Cervarix) reduce the risk of cancerous or precancerous changes of the cervix and perineum by about 93% and 62%, respectively. All have excellent safety profiles and are highly efficacious, or have met immunobridging standards. The vaccines are between 92% and 100% effective against HPV 16 and 18 up to at least 8 years.
HPV vaccines are typically given to age 9 to 26, as the vaccine is most effective if given before infection occurs. The primary target group in most of the countries recommending HPV vaccination is young adolescent girls, aged 9-14. As of 2022, 125 countries include HPV vaccines in their routine vaccinations for girls, and 47 countries recommend them for boys, as well, including Japan.
The duration of effectiveness and whether a booster will be needed is unknown. The high cost of this vaccine has been a cause for concern. Several countries have considered (or are considering) programs to fund HPV vaccination. The American Society of Clinical Oncology guideline has recommendations for different levels of resource availability.
Barrier protection
Barrier protection or spermicidal gel use during sexual intercourse decreases, but does not eliminate risk of transmitting the infection. Condoms may protect against genital warts. They also provide protection against other sexually transmitted infections, such as HIV and Chlamydia, which are associated with greater risks of developing cervical cancer.
Nutrition
Vitamin A is associated with a lower risk as are vitamin B12, vitamin C, vitamin E, and beta-Carotene.
Treatment
The treatment of cervical cancer varies worldwide, largely due to access to surgeons skilled in radical pelvic surgery and the emergence of fertility-sparing therapy in developed nations. Less advanced stages of cervical cancer typically have treatment options that allow fertility to be maintained if the patient desires.
Because cervical cancers are radiosensitive, radiation may be used in all stages where surgical options do not exist. Surgical intervention may have better outcomes than radiological approaches. In addition, chemotherapy can be used to treat cervical cancer and is more effective than radiation alone. Chemoradiotherapy may increase overall survival and reduce the risk of disease recurrence compared to radiotherapy alone.
Precancerous cells (cervical intraepithelial neoplasia) that would lead to cancer and early-stage cervical cancer (IA1) can be treated effectively by various surgical techniques. Surgical treatment methods include excision, where a cone-shaped portion of the cervix is removed, and ablation which removes only the parts with abnormal tissues. While these effectively reduce the risk of cancer developing or spreading, they cause an increased risk of premature birth in future pregnancies. Surgical techniques that remove more cervical tissue come with less risk of the cancer recurring but a higher chance of giving birth prematurely. Due to this risk, taking into account the age, childbearing plans of the woman, the size and location of the cancer cells are crucial for choosing the right procedure. There is low-certainty evidence that peri-operative care approaches, such as 'fast-track surgery' or 'enhanced recovery programmes' may lower surgical stress and improve recovery after gynaecological cancer surgery.
Microinvasive cancer (stage IA) may also be treated by hysterectomy (removal of the whole uterus, including part of the vagina). For stage IA2, the lymph nodes are removed as well. Alternatives include local surgical procedures such as a loop electrical excision procedure or cone biopsy. A systematic review concluded that more evidence is needed to inform decisions about different surgical techniques for women with cervical cancer at stage IA2.
If a cone biopsy does not produce clear margins (findings on biopsy showing that the tumor is surrounded by cancer free tissue, suggesting all of the tumor is removed), one more possible treatment option for women who want to preserve their fertility is a trachelectomy. This attempts to surgically remove the cancer while preserving the ovaries and uterus, providing for a more conservative operation than a hysterectomy. It is a viable option for those in stage I cervical cancer which has not spread; however, it is not yet considered a standard of care, as few doctors are skilled in this procedure. Even the most experienced surgeon cannot promise that a trachelectomy can be performed until after surgical microscopic examination, as the extent of the spread of cancer is unknown. If the surgeon is not able to microscopically confirm clear margins of cervical tissue once the woman is under general anaesthesia in the operating room, a hysterectomy may still be needed. This can only be done during the same operation if the woman consented. Due to the possible risk of cancer spreading to the lymph nodes in stage 1B cancers and some stage 1A cancers, the surgeon may also need to remove some lymph nodes from around the uterus for pathologic evaluation.
A radical trachelectomy can be performed abdominally or vaginally and opinions are conflicting as to which is better. A radical abdominal trachelectomy with lymphadenectomy usually only requires a two- to three-day hospital stay, and most women recover very quickly (about six weeks). Complications are uncommon, although women who can conceive after surgery are susceptible to preterm labour and possible late miscarriage. A wait of at least one year is generally recommended before attempting to become pregnant after surgery. Recurrence in the residual cervix is rare if the trachelectomy has cleared the cancer. Yet, women are recommended to practice vigilant prevention and follow-up care, including Pap screenings/colposcopy, with biopsies of the remaining lower uterine segment as needed (every 3–4 months for at least 5 years) to monitor for any recurrence in addition to minimizing any new exposures to HPV through safe sex practices until one is actively trying to conceive.
Early stages (IB1 and IIA less than 4 cm) can be treated with radical hysterectomy with removal of the lymph nodes or radiation therapy. Radiation therapy is given as external beam radiotherapy to the pelvis and brachytherapy (internal radiation). Women treated with surgery who have high-risk features found on pathologic examination are given radiation therapy with or without chemotherapy to reduce the risk of relapse. A Cochrane review has found moderate-certainty evidence that radiation decreases the risk of disease progression in people with stage IB cervical cancer, when compared to no further treatment. However, little evidence was found on its effects on overall survival.
Larger early-stage tumors (IB2 and IIA more than 4 cm) may be treated with radiation therapy and cisplatin-based chemotherapy, hysterectomy (which then usually requires adjuvant radiation therapy), or cisplatin chemotherapy followed by hysterectomy. When cisplatin is present, it is thought to be the most active single agent in periodic diseases. Such addition of platinum-based chemotherapy to chemoradiation seems not only to improve survival but also reduces risk of recurrence in women with early stage cervical cancer (IA2–IIA). A Cochrane review found a lack of evidence on the benefits and harms of primary hysterectomy compared to primary chemoradiotherapy for cervical cancer in stage IB2.
Advanced-stage tumors (IIB-IVA) are treated with radiation therapy and cisplatin-based chemotherapy. On 15 June 2006, the US Food and Drug Administration approved the use of a combination of two chemotherapy drugs, hycamtin and cisplatin, for women with late-stage (IVB) cervical cancer treatment. Combination treatment has significant risk of neutropenia, anemia, and thrombocytopenia side effects.
There is insufficient evidence whether anticancer drugs after standard care help women with locally advanced cervical cancer to live longer.
For surgery to be curative, the entire cancer must be removed with no cancer found at the margins of the removed tissue on examination under a microscope. This procedure is known as exenteration.
No evidence is available to suggest that any form of follow-up approach is better or worse in terms of prolonging survival, improving quality of life or guiding the management of problems that can arise because of the treatment and that in the case of radiotherapy treatment worsen with time. A 2019 review found no controlled trials regarding the efficacy and safety of interventions for vaginal bleeding in women with advanced cervical cancer.
Immunotherapy with immune checkpoint inhibitors, such as pembrolizumab (Keytruda), has also been approved by the U.S. Food and Drug Administration (FDA) for certain patients with recurrent or metastatic cervical cancer, demonstrating promising results in ongoing clinical trials. In October 2021, the FDA expanded this approval to include pembrolizumab in combination with chemotherapy, with or without bevacizumab, for people with persistent, recurrent, or metastatic cervical cancer, underscoring the potential of immunotherapeutic approaches in this setting. Additional immunotherapy agents, including other PD-1 and PD-L1 inhibitors, are under investigation and have similarly shown encouraging outcomes in clinical studies.
Another immune checkpoint inhibitor, cemiplimab-rwlc (Libtayo), received FDA approval in September 2022 for patients with recurrent or metastatic cervical cancer that has progressed on or after chemotherapy, further highlighting the expanding role of immunotherapeutic strategies in advanced disease.
Tisotumab vedotin (Tivdak) was approved for medical use in the United States in September 2021.
Prognosis
Stage
Prognosis depends on the stage of the cancer. For intraepithelial cervical neoplasmas the prognosis is good. With treatment, the five-year relative survival rate for the earliest stage of invasive cervical cancer is 92%, and the overall (all stages combined) five-year survival rate is about 66%. These statistics may be improved when applied to women newly diagnosed, bearing in mind that these outcomes may be partly based on the state of treatment five years ago when the women studied were first diagnosed.
With treatment, 80–90% of women with stage I cancer and 60–75% of those with stage II cancer are alive 5 years after diagnosis. Survival rates decrease to 58% for women with stage III cancer and 17% or fewer of those with stage IV cancer five years after diagnosis. Recurrent cervical cancer detected at its earliest stages might be successfully treated with surgery, radiation, chemotherapy, or a combination of the three. About 35% of women with invasive cervical cancer have persistent or recurrent disease after treatment.
By country
There is an ethnic disparity in five-year survival in the United States. Average survival rates of the dominate squamous cell carcinoma are 72% for Hispanic and Asian-Pacific women, 68% for White women and 61% for Black women.
Regular screening has meant that precancerous changes and early-stage cervical cancers have been detected and treated early. Figures suggest that cervical screening is saving 5,000 lives each year in the UK by preventing cervical cancer.
About 1,000 women per year die of cervical cancer in the UK. All of the Nordic countries have cervical cancer-screening programs in place. The Pap test was integrated into clinical practice in the Nordic countries in the 1960s.
In Africa outcomes are often worse as diagnosis is frequently at a latter stage of disease. In a scoping review of publicly-available cervical cancer prevention and control plans from African countries, plans tended to emphasize survivorship rather than early HPV diagnosis and prevention.
Adverse Effects
Chemotherapy works by attacking cells that rapidly divide. This kills cancer cells, but can also impact normal cells leading to adverse side effects. Common chemotherapy side effects include; hair loss, mouth sores, loss of appetite, diarrhea, nausea and vomiting, premature menopause, infertility, and damage to the blood-forming cells within bone marrow. Most acute side effects are temporary, dissipating when treatment ceases, but some can be long-lasting or permanent. Long-term chemotherapy side effects include changes in the menstrual cycle, neuropathy, and nephrotoxicity.
Radiation therapy (RT) adverse effects; for a complete side effect list see
Curative cervical radiation therapy may affect unintended tissues located within the delivery pathway(s) or adjacent to the target lesion, each tissue with a unique sensitivity and response to radiation injury. Common acute RT effects involve the gastrointestinal system, e.g., diarrhea and constipation; urinary tract, e.g., frequent urination; and may cause cervicitis. Common late RT complications include: infertility or premature ovarian failure; vaginal stenosis; lower motor neuron syndrome; telangiectasias, and subsequent hemorrhage; and progressive myelopathy, which may result in irreversible neurologic deficits ranging from minor sensory symptoms to complete paraplegia. Radiotherapy late effects (with occurrence rates) include osteonecrosis (8-20%), bladder ulceration (<3%), vaginal stenosis (>2.5%) and chronic pelvic radiation disease (1-10%), e.g., irreversible lumbosacral plexopathy.
Pelvic radiation also induces secondary malignancies such as leukemia, lymphoma, bladder cancer, pelvic malignancy, colorectal cancer, bone and soft-tissue sarcoma with occurrence rates between 0.2-1.0% per year for each.
Epidemiology
Worldwide, cervical cancer is both the fourth-most common type of cancer and the fourth-most common cause of death from cancer in women, with over 660,000 new cases and around 350,000 deaths in 2022. It is the second-most common cause of female-specific cancer after breast cancer, accounting for around 8% of both total cancer cases and total cancer deaths in women. 88% (2020 figure) of cervical cancers and 90% of deaths occur in low- and middle-income countries (LMICs) and 2% (2020 figure) in high-income countries (HICs). It is the most frequently detected cancer during pregnancy, with an occurrence of 1.5 to 12 for every 100,000 pregnancies.
The large majority of cervical cancer cases in 2020 (88%) occurred in LMICs, where they account for 17% of all cancers in women, compared with only 2% in high-income countries (HICs). In sub-Saharan Africa, the region with the highest rates of young women living with HIV (WLWH), approximately 20% of cervical cancer cases occur in WLWH. HPV infection is more likely to persist and to progress to cancer in WLWH. Mortality rates vary 50-fold between countries, ranging from <2 per 100 000 women in some HICs to >40 per 100 000 in some countries of sub-Saharan Africa.
Africa
Of the 20 hardest hit countries by cervical cancer, 19 are in Africa.
Australia
Australia is on target to eliminate cervical cancer. It anticipates to achieve this in the next 10 years.
In 2022, it is estimated that 942 new cases of cervical cancer will be diagnosed in Australia. In 2022, it is estimated that a female has a 1 in 180 (or 0.56%) risk of being diagnosed with cervical cancer by the age of 85.
In 2020, there were 165 women aged 25–74 who died from cervical cancer, which is a mortality rate of 2 deaths per 100,000 women in the population. Over the 5 years 2016–2020, there were 62 Aboriginal and Torres Strait Islander women aged 25–74 who died from cervical cancer, which is a mortality rate of 7 deaths per 100,000 Indigenous women in the population. Over the 5 years 2016–2020, the age-standardised mortality rate among Aboriginal and Torres Strait Islander women was 3.8 times the rate of non-Indigenous Australians.
The number of women diagnosed with cervical cancer has dropped on average by 4.5% each year since organised screening began in 1991 (1991–2005). Regular twice-yearly Pap tests can reduce the incidence of cervical cancer up to 90% in Australia, and save 1,200 Australian women from dying from the disease each year. It is predicted that because of the success of the primary HPV testing programme there will be fewer than four new cases per 100 000 women annually by 2028.
Canada
An estimated 1,450 Canadians will be diagnosed with cervical cancer in 2022. An estimated 380 will die from it.
India
In India, the number of people with cervical cancer is rising, but overall the age-adjusted rates are decreasing. Usage of condoms in the female population has improved the survival of women with cancers of the cervix.
European Union
As of 2022, the World Health Organization announced that "each year in the WHO European Region more than 66 000 women are newly diagnosed with cervical cancer and more than 30 000 die from this preventable disease."
United Kingdom
Cervical cancer is the 12th-most common cancer in women in the UK (around 3,100 women were diagnosed with the disease in 2011), and accounts for 1% of cancer deaths (around 920 died in 2012). With a 42% reduction from 1988 to 1997, the NHS-implemented screening programme has been highly successful, screening the highest-risk age group (25–49 years) every 3 years, and those ages 50–64 every 5 years.
United States
An estimated 13,170 new cervical cancers and 4,250 cervical cancer deaths will occur in the United States in 2019. The median age at diagnosis is 50. The rates of new cases in the United States was 7.3 per 100,000 women, based on rates from 2012 to 2016. Cervical cancer deaths decreased by approximately 74% in the last 50 years, largely due to widespread Pap test screening. The annual direct medical cost of cervical cancer prevention and treatment prior to introduction of the HPV vaccine was estimated at $6 billion.
Nigeria
The Nigerian Institute of Medical Research (NIMR) reports that 28 Nigerian women lose their lives daily due to this disease. This alarming statistic underscores the pressing need for better awareness, prevention, and treatment efforts across the country. Numerous Nigerian women lack access to these preventive measures. In many regions of the country, screening tests such as Pap tests and HPV tests are not easily accessible or affordable
History
400 BCE: Hippocrates noted that cervical cancer was incurable.
1925: Hinselmann invented the colposcope.
1928: Papanicolaou developed the Papanicolaou technique.
1941: Papanicolaou and Traut: Pap test screening began.
1946: Aylesbury spatula was developed to scrape the cervix, collecting the sample for the Pap test.
1951: First successful in-vitro cell line, HeLa, derived from biopsy of cervical cancer of Henrietta Lacks.
1976: Harald zur Hausen and Gisam found HPV DNA in cervical cancer and genital warts; Hausen later won the Nobel Prize for his work.
1988: Bethesda System for reporting Pap results was developed.
2006: First HPV vaccine was approved by the FDA.
2015: HPV vaccine shown to protect against infection at multiple body sites.
2018: Evidence for single-dose protection with HPV vaccine.
Epidemiologists working in the early 20th century noted that cervical cancer behaved like a sexually transmitted disease. In summary:
Cervical cancer was noted to be common in female sex workers.
It was rare in nuns, except for those who had been sexually active before entering the convent (Rigoni in 1841).
It was more common in the second wives of men whose first wives had died from cervical cancer.
It was rare in Jewish women.
In 1935, Syverton and Berry discovered a relationship between RPV (Rabbit Papillomavirus) and skin cancer in rabbits. (HPV is species-specific and therefore cannot be transmitted to rabbits).
These historical observations suggested that cervical cancer could be caused by a sexually transmitted agent. Initial research in the 1940s and 1950s attributed cervical cancer to smegma. During the 1960s and 1970s it was suspected that infection with herpes simplex virus (HSV) was the cause of the disease. In summary, HSV was seen as a likely cause because it is known to survive in the female reproductive tract, to be transmitted sexually in a way compatible with known risk factors, such as promiscuity and low socioeconomic status. Herpes viruses were also implicated in other malignant diseases, including Burkitt's lymphoma, Nasopharyngeal carcinoma, Marek's disease and the Lucké renal adenocarcinoma. HSV was recovered from cervical tumour cells.
A description of human papillomavirus (HPV) by electron microscopy was given in 1949, and HPV-DNA was identified in 1963. It was not until the 1980s that HPV was identified in cervical cancer tissue. It has since been demonstrated that HPV is implicated in virtually all cervical cancers. Specific viral subtypes implicated are HPV 16, 18, 31, 45 and others.
In work that was initiated in the mid-1980s, the HPV vaccine was developed, in parallel, by researchers at Georgetown University Medical Center, the University of Rochester, the University of Queensland in Australia, and the U.S. National Cancer Institute. In 2006, the US Food and Drug Administration (FDA) approved the first preventive HPV vaccine, marketed by Merck & Co. under the trade name Gardasil.
17 November is the Cervical Cancer Elimination Day of Action. The date marks the day in 2020 when WHO launched the Global strategy to accelerate the elimination of cervical cancer as a public health problem, with a resolution passed by 194 countries. In November 2020, the World Health Organization (WHO), under backing from the World Health Assembly, set out a strategy to eliminate cervical cancer by 2050. The strategy involves vaccinating 90% of girls by the age of 15, screening 70% of women by the age of 35 and again by the age of 45, and treating 90% of women identified with cervical disease. To eliminate cervical cancer, all countries must reach and maintain an incidence rate of below 4 per 100 000 women.
Society and culture
Australia
In Australia, Aboriginal women are more than five times more likely to die from cervical cancer than non-Aboriginal women, suggesting that Aboriginal women are less likely to have regular Pap tests. There are several factors that may limit indigenous women from engaging in regular cervical screening practices, including sensitivity in discussing the topic in Aboriginal communities, embarrassment, anxiety and fear about the procedure. Difficulty in accessing screening services (for example, transport difficulties) and a lack of female GPs, trained Pap test providers and trained female Aboriginal Health Workers are also issues.
The Australian Cervical Cancer Foundation (ACCF), founded in 2008, promotes 'women's health by eliminating cervical cancer and enabling treatment for women with cervical cancer and related health issues, in Australia and in developing countries.' Ian Frazer, one of the developers of the Gardasil cervical cancer vaccine, is the scientific advisor to ACCF. Janette Howard, the wife of the then-Prime Minister of Australia, John Howard, was diagnosed with cervical cancer in 1996, and first spoke publicly about the disease in 2006.
United States
A 2007 survey of American women found 40% had heard of HPV infection and less than half of those knew it causes cervical cancer. Over a longitudinal study from 1975 to 2000, it was found that people of lower socioeconomic census brackets had higher rates of late-stage cancer diagnosis and higher morbidity rates. After controlling for stage, there still existed differences in survival rates. Women in the US experience stigma around HPV infection, vaccination and cervical cancer. This is predominantly driven by fear of social judgment and rejection, self-blame, and shame, with notable negative influences from gender and social norms, as both human papillomavirus infection and cervical cancer were stigmatized due to the perception that they arise from reckless behavior such as having multiple sexual partners or neglecting screening.
LGBTQ populations
Transgender men and gender-diverse people who have a cervix (even if partially intact) or have a prior history of cervical cancer or precancerous conditions, and are age 21 or older who has ever had sex with anyone need to get screened for cervical cancer. Transmasculine people are just as likely as cisgender women to have cervical cancer, but are less likely to undergo cervical screening, because of dysphoria, gender disaffirmation or disempowerment of the individual by healthcare providers, or being misinformed of HPV and cervical cancer risks as well as many healthcare providers perceiving transmasculine individuals to be at low risk of cervical cancer.
Transgender women who have not had bottom surgery have no risk of cervical cancer, as they do not have a cervix. Trans women who have had bottom surgery to create a vagina (vaginoplasty) and possibly a cervix, are at a very small risk to develop cancer in the tissues of their neo-vagina or neo-cervix as these tissues are made up of different cells than a cervix in a cisgender woman Cervical cancer screening is not necessary in trans women who have undergone vaginoplasty because they do not have a cervix.
Intersex people with a cervix are also able to have cervical cancer.
| Biology and health sciences | Cancer | Health |
53436 | https://en.wikipedia.org/wiki/Bayonet | Bayonet | A bayonet (from Old French , now spelt ) is a knife, dagger, sword, or spike-shaped melee weapon designed to be mounted on the end of the barrel of a rifle, carbine, musket or similar long firearm, allowing the gun to be used as an improvised spear in close combat.
The term is derived from the town of Bayonne in southwestern France, where bayonets were supposedly first used by Basques in the 17th century. From the early 17th to the early 20th century, it was an infantry melee weapon used for both offensive and defensive tactics, usually when charging in mass formations (human wave attacks). In contemporary times, bayonets are considered a weapon of last resort, and are rarely used in combat, although they are still used for ceremonial purposes (e.g, Military parades).
History
The term bayonette itself dates back to the 16th century, but it is not clear whether bayonets at the time were knives that could be fitted to the ends of firearms, or simply a type of knife. For example, Cotgrave's 1611 Dictionarie describes the bayonet as "a kind of small flat pocket dagger, furnished with knives; or a great knife to hang at the girdle". Likewise, Pierre Borel wrote in 1655 that a kind of long-knife called a bayonette was made in Bayonne but does not give any further description. There are some accounts that place the invention of the bayonet in either France or Germany as early as 1570.
Plug bayonets
The first recorded instance of a bayonet proper is found in the Chinese military treatise, published in 1606. It was in the form of the , a breech-loading musket that was issued with a roughly long plug bayonet, giving it an overall length of with the bayonet attached. It was labelled as a "gun-blade" () with it being described as a "short sword that can be inserted into the barrel and secured by twisting it slightly" that it is to be used "when the battle have depleted both gunpowder and bullets as well as fighting against bandits, when forces are closing into melee or encountering an ambush" and if one "cannot load the gun within the time it takes to cover two bu (3.2 meters) of ground they are to attach the bayonet and hold it like a spear".
Early bayonets were of the "plug" type, where the bayonet was fitted directly into the barrel of the musket. This allowed light infantry to be converted to heavy infantry and hold off cavalry charges. The bayonet had a round handle that slid directly into the musket barrel. This naturally prevented the gun from being fired. The first known mention of the use of bayonets in European warfare was in the memoirs of Jacques de Chastenet, Vicomte de Puységur. He described the French using crude plug bayonets during the Thirty Years' War (1618–1648). However, it was not until 1671 that General Jean Martinet standardized and issued plug bayonets to the French regiment of fusiliers then raised. They were issued to part of an English dragoon regiment raised in 1672, and to the Royal Fusiliers when raised in 1685.
Socket bayonets
The major problem with plug bayonets was that when attached they made it impossible to fire the musket, requiring soldiers to wait until the last possible moment before a melee to fix the bayonet. The defeat of forces loyal to William of Orange by Jacobite Highlanders at the Battle of Killiecrankie in 1689 was due (among other things) to the use of the plug bayonet. The Highlanders closed to , fired a single volley, dropped their muskets, and using axes and swords quickly overwhelmed the loyalists before they had time to fix bayonets. Shortly thereafter, the defeated leader, Hugh Mackay, is believed to have introduced a socket bayonet of his own invention. Soon "socket" bayonets would incorporate both socket mounts and an offset blade that fit around the musket's barrel, which allowed the musket to be fired and reloaded while the bayonet was attached.
An unsuccessful trial with socket or zigzag bayonets was made after the Battle of Fleurus in 1690, in the presence of King Louis XIV, who refused to adopt them, as they had a tendency to fall off the musket. Shortly after the Peace of Ryswick (1697), the English and Germans abolished the pike and introduced socket bayonets. The British socket bayonet had a spike with a triangular cross-section rather than a flat blade, with a flat side towards the muzzle and two fluted sides outermost to a length of . It had no lock to keep it fast to the muzzle, and was well-documented for falling off in the heat of battle.
By the mid-18th century, socket bayonets had been adopted by most European armies. In 1703, the French infantry adopted a spring-loaded locking system that prevented the bayonet from accidentally separating from the musket. A triangular blade was introduced around 1715 and was stronger than the previous single or double-edged model.
Sword bayonets
The 18th century introduced the concept of the sword bayonet, a long-bladed weapon with a single- or double-edged blade that could also be used as a shortsword. Its initial purpose was to ensure that riflemen could form an infantry square properly to fend off cavalry attacks when in ranks with musketmen, whose weapons were longer. A prime early example of a sword bayonet-fitted rifle is the Pattern 1800 Infantry Rifle, later known as the "Baker Rifle". Sword bayonets were used by German Jagers in the 18th century. The hilt usually had quillons modified to accommodate the gun barrel and a hilt mechanism that enabled the bayonet to be attached to a bayonet lug. A sword bayonet could be used in combat as a sidearm, when detached from the musket or rifle. When the bayonet was attached to the musket or rifle, it effectively turned all long guns into a spear or glaive, which made it suitable for both thrusting and cutting attacks.
While the British Army eventually discarded the sword bayonet, the socket bayonet survived the introduction of the rifled musket into British service in 1854. The new rifled musket copied the French locking ring system. The new bayonet proved its worth at the Battle of Alma and the Battle of Inkerman during the Crimean War, where the Imperial Russian Army learned to fear it.
In the 1860s, European nations began to develop new bolt-action breechloading rifles (such as the Chassepot and Snider–Enfield) and sword bayonets suitable for mass production and used by police, pioneer, and engineer troops. The decision to redesign the bayonet into a short sword was viewed by some as an acknowledgement of the decline in importance of the fixed bayonet as a weapon in the face of new advances in firearms technology. As a British newspaper put it, "the committee, in recommending this new sword bayonet, appear to have had in view the fact that bayonets will henceforth be less frequently used than in former times as a weapon of offence and defence; they desired, therefore, to substitute an instrument of more general utility."
Multipurpose bayonets
One of these multipurpose designs was the 'sawback' bayonet, which incorporated saw teeth on the spine of the blade. The sawback bayonet was intended for use as a general-purpose utility tool as well as a weapon; the teeth were meant to facilitate the cutting of wood for various defensive works such as barbed-wire posts, as well as for butchering livestock. It was initially adopted by the German states in 1865; until the middle of WWI approximately 5% of every bayonet style was complemented with a sawback version, for example in Belgium in 1868, Great Britain in 1869 and Switzerland in 1878 (Switzerland introduced their last model in 1914). The original sawback bayonets were typically of the heavy sword-type, they were issued to engineers, with to some extent the bayonet aspect being secondary to the "tool" aspect. Later German sawbacks were more of a rank indicator than a functional saw. Generally, an average of 6% of all bayonets were sawbacks for non-commissioned officers. There were some exceptions, such as the kurzes Seitengewehr 1898 model, all of which were of the sawback design and meant for what was considered more prestigious units, such as machine gunners, telegraph troop and colonial troops. The sawback proved relatively ineffective as a cutting tool, and was soon outmoded by improvements in military logistics and transportation; most nations dropped the sawback feature by the early 20th century. The German army discontinued use of the sawback bayonet in 1917 after protests that the serrated blade caused unnecessarily severe wounds when used as a fixed bayonet.
The trowel or spade bayonet was another multipurpose design, intended for use both as an offensive weapon as well as a digging tool for excavating entrenchments. In 1870, the US Army issued trowel bayonets to infantry regiments based on a design by Lieutenant-Colonel Edmund Rice, a US Army officer and Civil War veteran, which were manufactured by the Springfield Armory. Besides its utility as both a fixed bayonet and a digging implement, the Rice trowel bayonet could be used to plaster log huts and stone chimneys for winter quarters; sharpened on one edge, it could cut tent poles and pins. Ten thousand were eventually issued, and the design saw service during the 1877 Nez Perce campaign. Rice was given leave in 1877 to demonstrate his trowel bayonet to several nations in Europe. One infantry officer recommended it to the exclusion of all other designs, noting that "the entrenching tools of an army rarely get up to the front until the exigency for their use has passed." The Rice trowel bayonet was declared obsolete by the US Army in December 1881.
"Reach" controversy
Prior to World War I, bayonet doctrine was largely founded upon the concept of "reach"; that is, a soldier's theoretical ability, by use of an extremely long rifle and fixed bayonet, to stab an enemy soldier without having to approach within reach of his opponent's blade. A combined length of rifle and bayonet longer than that of the enemy infantryman's rifle and attached bayonet, like the infantryman's pike of bygone days, was thought to impart a tactical advantage on the battlefield.
In 1886, the French army introduced a quadrangular épée spike for the bayonet of the Lebel Model 1886 rifle, the Épée-Baïonnette Modèle 1886, resulting in a rifle and bayonet with an overall length of . Germany responded by introducing a long sword bayonet for the Model 1898 Mauser rifle, which had a 29-inch barrel. The bayonet, the Seitengewehr 98, had a 50 cm (19.7-inch) blade. With an overall length of , the German army's rifle/bayonet combination was second only to the French Lebel for overall 'reach'.
After 1900, Switzerland, Britain, and the United States adopted rifles with barrel lengths shorter than that of a rifled musket, but longer than that of a carbine. These were intended for general use by infantry and cavalry. The "reach" of the new short rifles with attached bayonets was reduced. Britain introduced the SMLE (Short, Magazine, Lee–Enfield), in 1904. The German M1898 Mauser rifle and attached sword bayonet was 20 cm (eight inches) longer than the SMLE and its P1903 bayonet, which used a twelve-inch (30 cm) blade. While the British P1903 and its similar predecessor, the P1888, was satisfactory in service, criticism soon arose regarding the shortened reach. One military writer of the day warned: "The German soldier has eight inches the better of the argument over the British soldier when it comes to crossing bayonets, and the extra eight inches easily turns the battle in favour of the longer, if both men are of equal skill."
In 1905, the German Army adopted a shortened bayonet, the Seitengewehr 98/05 for engineer and pioneer troops, and in 1908, a short rifle as well, the Karabiner Model 1898AZ, which was produced in limited quantities for the cavalry, artillery, and other specialist troops. However, the long-barreled 98 Mauser rifle remained in service as the primary infantry small arm. Moreover, German military authorities continued to promote the idea of outreaching one's opponent on the battlefield by means of a longer rifle/bayonet combination, a concept prominently featured in its infantry bayonet training doctrines. These included the throw point or extended thrust-and-lunge attack. Using this tactic, the German soldier dropped into a half-crouch, with the rifle and fixed bayonet held close to the body. In this position the soldier next propelled his rifle forward, then dropped the supporting hand while taking a step forward with the right foot, simultaneously thrusting out the right arm to full length with the extended rifle held in the grip of the right hand alone. With a maximum 'kill zone' of some eleven feet, the throw point bayonet attack gave an impressive increase in 'reach', and was later adopted by other military forces, including the U.S. Army.
In response to criticism over the reduced reach of the SMLE rifle and bayonet, British ordnance authorities introduced the P1907 bayonet in 1908, which had an elongated blade of some seventeen inches to compensate for the reduced overall length of the SMLE rifle. The 1907 bayonet was essentially a copy of the Japanese Type 30 bayonet, Britain having purchased a number of Japanese type 30 rifles for the Royal Navy during the preceding years. U.S. authorities in turn adopted a long (16-in. blade) bayonet for the M1903 Springfield short rifle, the M1905 bayonet; later, a long sword bayonet was also provided for the M1917 Enfield rifle.
Reversal in opinion
The experience of World War I reversed opinion on the value of long rifles and bayonets in typical infantry combat operations. Whether in the close confines of trench warfare, night time raiding and patrolling, or attacking across open ground, soldiers of both sides soon recognized the inherent limitations of a long and ungainly rifle and bayonet when used as a close-quarters battle weapon. Once Allied soldiers had been trained to expect the throw point or extended thrust-and-lunge attack, the method lost most of its tactical value on the World War I battlefield. It required a strong arm and wrist, was very slow to recover if the initial thrust missed its mark, and was easily parried by a soldier who was trained to expect it, thus exposing the German soldier to a return thrust which he could not easily block or parry. Instead of longer bayonets, infantry forces on both sides began experimenting with other weapons as auxiliary close-quarter arms, including the trench knife, trench club, handgun, hand grenade, and entrenching tool.
Soldiers soon began employing the bayonet as a knife as well as an attachment for the rifle, and bayonets were often shortened officially or unofficially to make them more versatile and easier to use as tools, or to maneuver in close quarters. During World War II, bayonets were further shortened into knife-sized weapons in order to give them additional utility as fighting or utility knives. The vast majority of modern bayonets introduced since World War II are of the knife bayonet type.
Bayonet charge
The development of the bayonet from the 17th century onwards led to the bayonet charge becoming the main infantry tactic throughout the 18th, 19th, and early 20th century. The British Army under Wolesley, the later Duke of Wellington, evolved its tactics to adopt the "Volley and Bayonet Charge" from the earlier "Highland Charge" tactic of Highland regiments under his command. These proto "fire and maneuver" tactics were first introduced to the British Army by the 42nd Highlanders (Black Watch) at Fontenoy in 1745 although, they had been used by their antecedents, (The Independent Highland Watch Companies) prior to that. As early as the 19th century, military scholars were already noting that most bayonet charges did not result in close combat. Instead, one side usually fled before actual bayonet fighting ensued. The act of fixing bayonets has been held to be primarily connected to morale, the making of a clear signal to friend and foe of a willingness to kill at close quarters.
The bayonet charge was above all a tool of shock. While charges were reasonably common in 18th and 19th century warfare, actual combat between formations with their bayonets was so rare as to be effectively nonexistent. Usually, a charge would only happen after a long exchange of gunfire, and one side would break and run before contact was actually made. Sir Charles Oman, nearing the end of his history of the Peninsular War (1807–1814) in which he had closely studied hundreds of battles and combats, only discovered a single example of, in his words, "one of the rarest things in the Peninsular War, a real hand-to-hand fight with the white weapon." Infantry melees were much more common in close country – towns, villages, earthworks and other terrain which reduced visibility to such ranges that hand-to-hand fighting was unavoidable. These melees, however, were not bayonet charges per se, as they were not executed or defended against by regular bodies of orderly infantry; rather, they were a chaotic series of individual combats where musket butts and fists were used alongside bayonets, swords, and polearms.
Napoleonic wars
The bayonet charge was a common tactic used during the Napoleonic wars. Despite its effectiveness, a bayonet charge did not necessarily cause substantial casualties through the use of the weapon itself. Detailed battle casualty lists from the 18th century showed that in many battles, less than 2% of all wounds treated were caused by bayonets. Antoine-Henri Jomini, a celebrated military author who served in numerous armies during the Napoleonic period, stated that the majority of bayonet charges in the open resulted with one side fleeing before any contact was made. Combat with bayonets did occur, but mostly on a small scale when units of opposing sides encountered each other in a confined environment, such as during the storming of fortifications or during ambush skirmishes in broken terrain. In an age of fire by massed volley, when compared to random unseen bullets, the threat of the bayonet was much more tangible and immediate – guaranteed to lead to a personal gruesome conclusion if both sides persisted. All this encouraged men to flee before the lines met. Thus, the bayonet was an immensely useful weapon for capturing ground from the enemy, despite seldom actually being used to inflict wounds.
American Civil War
During the American Civil War (1861–1865) the bayonet was found to be responsible for less than 1% of battlefield casualties, a hallmark of modern warfare. The use of bayonet charges to force the enemy to retreat was very successful in numerous small unit engagements at short range in the American Civil War, as most troops would retreat when charged while reloading. Although such charges inflicted few casualties, they often decided short engagements, and tactical possession of important defensive ground features. Additionally, bayonet drill could be used to rally men temporarily unnerved by enemy fire.
While the overall Battle of Gettysburg was won by the Union armies due to a combination of terrain and massed artillery fire, a decisive point on the second day of the battle hinged on a bayonet charge at Little Round Top when Joshua Lawrence Chamberlain's 20th Maine Volunteer Infantry Regiment, running short of musket ammunition, charged downhill, surprising and capturing many of the surviving soldiers of the 15th Alabama Infantry Regiment and other Confederate regiments. Other bayonet charges occurred at Gettysburg, such as that of the 1st Minnesota Infantry Regiment. This was ordered in desperation by General Hancock earlier on July 2 in order to delay a Confederate brigade's advance long enough to bring up reinforcements for the holed Union line on Cemetery Ridge. Still another bayonet charge was conducted late in the evening on July 2 by the 137th New York Infantry Regiment defending the extreme right flank of the Union line on Culp's Hill. The charge of several companies managed to temporarily stall the advance of the 10th Virginia Infantry Regiment long enough for the 14th Brooklyn to move in on the 137th's right and repel the attack.
Going over the top
The popular image of World War I combat is of a wave of soldiers with bayonets fixed, "going over the top" and charging across no man's land into a hail of enemy fire. Although this was the standard method of fighting early in the war, it was rarely successful. British casualties on the first day of the Battle of the Somme were the worst in the history of the British army, with casualties, whom were killed.
During World War I, no man's land was often hundreds of yards across. The area was usually devastated by the warfare and riddled with craters from artillery and mortar shells, and sometimes contaminated by chemical weapons. Heavily defended by machine guns, mortars, artillery, and riflemen on both sides, it was often covered with barbed wire and land mines, and littered with the rotting corpses of those who were not able to make it across the sea of projectiles, explosions, and flames. A bayonet charge through no man's land often resulted in the total annihilation of entire battalions.
Banzai charges
The advent of modern warfare in the 20th century made bayonet charges dubious affairs. During the Siege of Port Arthur (1904–1905), the Japanese used human wave attacks against Russian artillery and machine guns, suffering massive casualties.
However, during the Second Sino-Japanese War, the Japanese were able to use bayonet charges effectively against poorly organized and lightly armed Chinese troops. "Banzai charges" became an accepted military tactic where Japanese forces were able to rout larger Chinese forces routinely.
In the early stages of the Pacific War (1941–1945), a sudden bayonet charge could overwhelm unprepared enemy soldiers. Such charges became known to Allied forces as "Banzai charges" from the Japanese battle cry. By the end of the war, against well organized and heavily armed Allied forces, a banzai charge inflicted little damage but at high cost. They were sometimes conducted as a last resort by small groups of surviving soldiers when the main battle was already lost.
Some Japanese commanders, such as General Tadamichi Kuribayashi, recognized the futility and waste of such attacks and expressly forbade their men from carrying them out. Indeed, the Americans were surprised that the Japanese did not employ banzai charges at the Battle of Iwo Jima.
Human wave attack
The term "human wave attack" was often misused to describe the Chinese short attack—a combination of infiltration and the shock tactics employed by the People's Liberation Army during the Korean War (1950–1953). A typical Chinese short attack was carried out at night by sending a series of small five-man fireteams to attack the weakest point of an enemy's defenses. The Chinese assault team would crawl undetected within grenade range, then launch surprise attacks with fixed bayonets against the defenders in order to breach the defenses by relying on maximum shock and confusion.
If the initial shock failed to breach the defenses, additional fireteams would press on behind them and attack the same point until a breach was created. Once penetration was achieved, the bulk of the Chinese forces would move into the enemy rear and attack from behind. Due to primitive communication systems and tight political controls within the Chinese army, short attacks were often repeated until either the defenses were penetrated or the attackers were completely annihilated.
This persistent attack pattern left a strong impression on UN forces that fought in Korea, giving birth to the description of "human wave". The term "human wave" was later used by journalists and military officials to convey the image of the American soldiers being assaulted by overwhelming numbers of Chinese on a broad front, which is inaccurate when compared with the normal Chinese practice of sending successive series of small teams against a weak point in the line. It was in fact rare for the Chinese to actually use densely concentrated infantry formations to absorb enemy firepower.
Modern usage
One use the Germans in World War II made of bayonets was to search for people in hiding. One person hiding in a house in the Netherlands wrote: "The Germans made lots of noise as they came upstairs, and they stabbed their bayonets into the wall. Then what we'd always feared actually happened: A bayonet went through the thin wallpaper above the closet, exposing the three people who were hiding there. 'Raus!' cried the Germans. 'Out!'".
During the Korean War, the French Battalion and Turkish Brigade used bayonet charges against enemy combatants. In 1951, United States Army officer Lewis L. Millett led soldiers of the US Army's 27th Infantry Regiment in capturing a machine gun position with bayonets. Historian S. L. A. Marshall described the attack as "the most complete bayonet charge by American troops since Cold Harbor". The location subsequently became known as Bayonet Hill. This was the last bayonet charge by the US Army. Millett was awarded the Medal of Honor.
On 23 October 1962, during the Sino-Indian War, 20 Indian soldiers led by Joginder Singh fixed bayonets and charged a force of 200 Chinese soldiers. While the charge would prove futile for Singh and his men, it initially threw the Chinese off guard and forced a retreat despite outnumbering them 10 to 1.
On 8 May 1970, National Guardsmen attacked student demonstrators with bayonets at the University of New Mexico in Albuquerque. The demonstrators were protesting the war in Vietnam and Cambodia, and the killing of four students at Kent State University. Eleven were injured, some seriously.
In 1982, the British Army mounted bayonet charges during the Falklands War, notably the 3rd Battalion, Parachute Regiment during the Battle of Mount Longdon and the 2nd Battalion, Scots Guards during the final assault of Mount Tumbledown.
In 1995, during the Siege of Sarajevo, UN peacekeepers of the French 3rd Marine Infantry Regiment charged Serbian forces at the Battle of Vrbanja bridge. Actions led by the regiment allowed the UN peacekeepers to retreat from a threatened position. Two fatalities and seventeen wounded resulted.
During the Second Gulf War and the war in Afghanistan, British Army units mounted several bayonet charges. In 2004, at the Battle of Danny Boy in Iraq, the Argyll and Sutherland Highlanders charged mortar positions of the Mahdi Army. The ensuing hand-to-hand fighting resulted in an estimate of over 40 insurgents killed and 35 bodies collected and nine prisoners. Sergeant Brian Wood, of the Princess of Wales's Royal Regiment, was awarded the Military Cross for his part in the battle.
In 2009, Lieutenant James Adamson of the Royal Regiment of Scotland was awarded the Military Cross for a bayonet charge while in Afghanistan. Adamson had run out of ammunition so he immediately charged a Taliban fighter with his bayonet. Lance Corporal Sean Jones of The Princess of Wales's Regiment was awarded the Military Cross for his role in a 2011 bayonet charge.
Contemporary bayonets
Today, the bayonet is rarely used in one-to-one combat. Despite its limitations, many modern assault rifles (including bullpup designs) retain a bayonet lug and the bayonet is issued by many armies. The bayonet is used for controlling prisoners, or as a weapon of last resort. In addition, some authorities have concluded that the bayonet serves as a useful training aid in building morale and increasing desired aggressiveness in troops.
Today's bayonets often double as multi-purpose utility knives, bottle openers or other tools. Issuing one modern multi-purpose bayonet/knife is also more cost effective than issuing separate specialty bayonets, and field/combat knives.
Soviet Union
The original AK-47 has an adequate but unremarkable bayonet. However, the AKM Type I bayonet (introduced in 1959) was an improvement of the original design. It has a Bowie style (clip-point) blade with saw-teeth along the spine, and can be used as a multi-purpose survival knife and wire-cutter when combined with its steel scabbard. The AK-74 bayonet 6Kh5 (introduced in 1983) represents a further refinement of the AKM bayonet. "It introduced a radical blade cross-section, that has a flat milled on one side near the edge and a corresponding flat milled on the opposite side near the false edge." The blade has a new spear point and an improved one-piece moulded plastic grip, making it a more effective fighting knife. It also has saw-teeth on the false edge and the usual hole for use as a wire-cutter. The wire cutting versions of the AK bayonets each have an electrically insulated handle and an electrically insulated part of the scabbard, so it can be used to cut an electrified wire.
United States
The American M16 rifle used the M7 bayonet which is based on earlier designs such as the M4, M5 and M6 models, all of which are direct descendants of the M3 Fighting Knife and have a spear-point blade with a half sharpened secondary edge. The newer M9 has a clip-point blade with saw-teeth along the spine, and can be used as a multi-purpose knife and wire-cutter when combined with its scabbard. It can even be used by troops to cut their way free through the relatively thin metal skin of a crashed helicopter or airplane. The current USMC OKC-3S bayonet bears a resemblance to the Marines' iconic Ka-Bar fighting knife with serrations near the handle.
People's Republic of China
The AK-47 was adopted by Communist China as the Type 56 assault rifle and includes an integral folding spike bayonet, similar to the SKS rifle. Some Type 56s may also use the AKM Type II bayonet. The latest Chinese rifle, the QBZ-95, has a multi-purpose knife bayonet similar to the US M9.
Belgium
The FN FAL has two types of bayonet. The first is a traditional spear point bayonet. The second is the Type C socket bayonet introduced in the 1960s. It has a hollow handle that fits over the muzzle and slots that lined up with those on the FALs 22 mm NATO-spec flash hider. Its spear-type blade is offset to the side of the handle to allow the bullet to pass beside the blade.
United Kingdom
The current British L3A1 socket bayonet is based on the FN FAL Type C socket bayonet with a clip-point blade. It has a hollow handle that fits over the SA80/L85 rifle's muzzle and slots that lined up with those on the flash eliminator. The blade is offset to the side of the handle to allow the bullet to pass beside the blade. It can also be used as a multi-purpose knife and wire-cutter when combined with its scabbard. The scabbard also has a sharpening stone and folding saw blade. The use of contemporary bayonets by the British army was noted during the Afghanistan war in 2004. Traditionally, bayonets are instead called swords in The Rifles.
Germany
The H&K G3 rifle uses two types of bayonets, both of which is mounted above the G3's rifle barrel. The first is the standard G3 bayonet which has a blade similar to the American M7 bayonet. The second is an Eickhorn KCB-70 type multi-purpose knife bayonet, featuring a clip-point with saw-back, a wire-cutter scabbard and a distinctive squared handgrip. For the H&K G36 there was little use of modified AKM type II knife bayonets from stocks of the former Nationale Volksarmee (National People's Army) of East Germany. The original muzzle-ring was cut away and a new, large diameter muzzle ring welded in place. The original leather belt hanger was replaced by a complex web and plastic belt hanger designed to fit the West German load bearing equipment.
Austria
The Steyr AUG uses two types of bayonet. The first and most common is an Eickhorn KCB-70 type multi-purpose bayonet with an M16 bayonet type interface. The second are the Glock Feldmesser 78 (Field Knife 78) and the Feldmesser 81 (Survival Knife 81), which can also be used as a bayonet, by engaging a socket in the pommel (covered by a plastic cap) into a bayonet adapter that can be fitted to the AUG rifle. These bayonets are noteworthy, as they were meant to be used primarily as field or survival knives and use as a bayonet was a secondary consideration. They can also be used as throwing knives and have a built-in bottle opener in the crossguard.
France
The French use a more traditional spear point bayonet with the current FAMAS bayonet which is nearly identical to that of the M1949/56 bayonet. The new French H&K 416F rifle uses the Eickhorn "SG 2000 WC-F", a multi-purpose combat knife/bayonet (similar to the KM2000) with a wire cutter. It weighs , is long with a half serrated blade for cutting through ropes. The synthetic handle and sheath have electrical insulation that protects up to 10,000 volts. The sheath also has a diamond blade sharpener.
Photo gallery
Linguistic impact
The push-twist motion of fastening the older type of spike bayonet has given a name to:
The "bayonet mount" used for various types of quick fastenings, such as camera lenses, also called a "bayonet connector" when used in electrical plugs.
Several connectors and contacts including the bayonet-fitting light bulb that is common in the UK (as opposed to the continental European screw-fitting type).
One type of connector for foil and sabre weapons used in modern fencing competitions is referred to as a "bayonet" connector.
In chess, an aggressive variation of the King's Indian Defence is known as the "Bayonet Attack".
The bayonet has become a symbol of military power. The term "at the point of a bayonet" refers to using military force or action to accomplish, maintain, or defend something (cf. Bayonet Constitution). Undertaking a task "with fixed bayonets" has this connotation of no room for compromise and is a phrase used particularly in politics.
Badges and insignias
The Australian Army 'Rising Sun' badge features a semicircle of bayonets. The Australian Army Infantry Combat Badge (ICB) takes the form of a vertically mounted Australian Army SLR (7.62mm self-loading rifle FN FAL) bayonet surrounded by an oval-shaped laurel wreath.
The US Army Combat Action Badge, awarded to personnel who have come under fire since 2001 and who are not eligible for the Combat Infantryman Badge (due to the fact that only Infantry personnel may be awarded the Combat Infantryman Badge), has a bayonet as its central motif. The shoulder sleeve insignia for the 10th Mountain Division in the US Army features crossed bayonets. The US Army's 173rd Airborne Brigade Combat Team's shoulder patch features a bayonet wrapped in a wing, symbolizing their airborne status. The brigade regularly deploys in task forces under the name "Bayonet".
The insignia of the British Army's School of Infantry is an SA80 bayonet against a red shield. It is worn as a Tactical recognition flash (TRF) by instructors at the Infantry Training Centre Catterick, the Infantry Battle School at Brecon and the Support Weapons School in Warminster. Fixed bayonets also feature on the cap badge and tactical recognition flash of the Small Arms School Corps.
The vocation tab collar insignia for the Singapore Armed Forces Infantry Formation utilizes two crossed bayonets. The bayonet is often used as a symbol of the Infantry in Singapore.
| Technology | Melee weapons | null |
53452 | https://en.wikipedia.org/wiki/Euler%27s%20totient%20function | Euler's totient function | In number theory, Euler's totient function counts the positive integers up to a given integer that are relatively prime to . It is written using the Greek letter phi as or , and may also be called Euler's phi function. In other words, it is the number of integers in the range for which the greatest common divisor is equal to 1. The integers of this form are sometimes referred to as totatives of .
For example, the totatives of are the six numbers 1, 2, 4, 5, 7 and 8. They are all relatively prime to 9, but the other three numbers in this range, 3, 6, and 9 are not, since and . Therefore, . As another example, since for the only integer in the range from 1 to is 1 itself, and .
Euler's totient function is a multiplicative function, meaning that if two numbers and are relatively prime, then .
This function gives the order of the multiplicative group of integers modulo (the group of units of the ring ). It is also used for defining the RSA encryption system.
History, terminology, and notation
Leonhard Euler introduced the function in 1763. However, he did not at that time choose any specific symbol to denote it. In a 1784 publication, Euler studied the function further, choosing the Greek letter to denote it: he wrote for "the multitude of numbers less than , and which have no common divisor with it". This definition varies from the current definition for the totient function at but is otherwise the same. The now-standard notation comes from Gauss's 1801 treatise Disquisitiones Arithmeticae, although Gauss did not use parentheses around the argument and wrote . Thus, it is often called Euler's phi function or simply the phi function.
In 1879, J. J. Sylvester coined the term totient for this function, so it is also referred to as Euler's totient function, the Euler totient, or Euler's totient. Jordan's totient is a generalization of Euler's.
The cototient of is defined as . It counts the number of positive integers less than or equal to that have at least one prime factor in common with .
Computing Euler's totient function
There are several formulae for computing .
Euler's product formula
It states
where the product is over the distinct prime numbers dividing . (For notation, see Arithmetical function.)
An equivalent formulation is
where is the prime factorization of (that is, are distinct prime numbers).
The proof of these formulae depends on two important facts.
Phi is a multiplicative function
This means that if , then . Proof outline: Let , , be the sets of positive integers which are coprime to and less than , , , respectively, so that , etc. Then there is a bijection between and by the Chinese remainder theorem.
Value of phi for a prime power argument
If is prime and , then
Proof: Since is a prime number, the only possible values of are , and the only way to have is if is a multiple of , that is, , and there are such multiples not greater than . Therefore, the other numbers are all relatively prime to .
Proof of Euler's product formula
The fundamental theorem of arithmetic states that if there is a unique expression where are prime numbers and each . (The case corresponds to the empty product.) Repeatedly using the multiplicative property of and the formula for gives
This gives both versions of Euler's product formula.
An alternative proof that does not require the multiplicative property instead uses the inclusion-exclusion principle applied to the set , excluding the sets of integers divisible by the prime divisors.
Example
In words: the distinct prime factors of 20 are 2 and 5; half of the twenty integers from 1 to 20 are divisible by 2, leaving ten; a fifth of those are divisible by 5, leaving eight numbers coprime to 20; these are: 1, 3, 7, 9, 11, 13, 17, 19.
The alternative formula uses only integers:
Fourier transform
The totient is the discrete Fourier transform of the gcd, evaluated at 1. Let
where for . Then
The real part of this formula is
For example, using and :Unlike the Euler product and the divisor sum formula, this one does not require knowing the factors of . However, it does involve the calculation of the greatest common divisor of and every positive integer less than , which suffices to provide the factorization anyway.
Divisor sum
The property established by Gauss, that
where the sum is over all positive divisors of , can be proven in several ways. (See Arithmetical function for notational conventions.)
One proof is to note that is also equal to the number of possible generators of the cyclic group ; specifically, if with , then is a generator for every coprime to . Since every element of generates a cyclic subgroup, and each subgroup is generated by precisely elements of , the formula follows. Equivalently, the formula can be derived by the same argument applied to the multiplicative group of the th roots of unity and the primitive th roots of unity.
The formula can also be derived from elementary arithmetic. For example, let and consider the positive fractions up to 1 with denominator 20:
Put them into lowest terms:
These twenty fractions are all the positive ≤ 1 whose denominators are the divisors . The fractions with 20 as denominator are those with numerators relatively prime to 20, namely , , , , , , , ; by definition this is fractions. Similarly, there are fractions with denominator 10, and fractions with denominator 5, etc. Thus the set of twenty fractions is split into subsets of size for each dividing 20. A similar argument applies for any n.
Möbius inversion applied to the divisor sum formula gives
where is the Möbius function, the multiplicative function defined by and for each prime and . This formula may also be derived from the product formula by multiplying out to get
An example:
Some values
The first 100 values are shown in the table and graph below:
{| class="wikitable" style="text-align: right"
|+ for
! +
! 1 || 2 || 3 || 4 || 5 || 6 || 7 || 8 || 9 || 10
|-
! 0
| 1 || 1 || 2 || 2 || 4 || 2 || 6 || 4 || 6 || 4
|-
! 10
| 10 || 4 || 12 || 6 || 8 || 8 || 16 || 6 || 18 || 8
|-
! 20
| 12 || 10 || 22 || 8 || 20 || 12 || 18 || 12 || 28 || 8
|-
! 30
| 30 || 16 || 20 || 16 || 24 || 12 || 36 || 18 || 24 || 16
|-
! 40
| 40 || 12 || 42 || 20 || 24 || 22 || 46 || 16 || 42 || 20
|-
! 50
| 32 || 24 || 52 || 18 || 40 || 24 || 36 || 28 || 58 || 16
|-
! 60
| 60 || 30 || 36 || 32 || 48 || 20 || 66 || 32 || 44 || 24
|-
! 70
| 70 || 24 || 72 || 36 || 40 || 36 || 60 || 24 || 78 || 32
|-
! 80
| 54 || 40 || 82 || 24 || 64 || 42 || 56 || 40 || 88 || 24
|-
! 90
| 72 || 44 || 60 || 46 || 72 || 32 || 96 || 42 || 60 || 40
|}
In the graph at right the top line is an upper bound valid for all other than one, and attained if and only if is a prime number. A simple lower bound is , which is rather loose: in fact, the lower limit of the graph is proportional to .
Euler's theorem
This states that if and are relatively prime then
The special case where is prime is known as Fermat's little theorem.
This follows from Lagrange's theorem and the fact that is the order of the multiplicative group of integers modulo .
The RSA cryptosystem is based on this theorem: it implies that the inverse of the function , where is the (public) encryption exponent, is the function , where , the (private) decryption exponent, is the multiplicative inverse of modulo . The difficulty of computing without knowing the factorization of is thus the difficulty of computing : this is known as the RSA problem which can be solved by factoring . The owner of the private key knows the factorization, since an RSA private key is constructed by choosing as the product of two (randomly chosen) large primes and . Only is publicly disclosed, and given the difficulty to factor large numbers we have the guarantee that no one else knows the factorization.
Other formulae
In particular:
Compare this to the formula (see least common multiple).
is even for . Moreover, if has distinct odd prime factors,
For any and such that there exists an such that .
where is the radical of (the product of all distinct primes dividing ).
( cited in)
[Liu (2016)]
(where is the Euler–Mascheroni constant).
Menon's identity
In 1965 P. Kesava Menon proved
where is the number of divisors of .
Divisibility by any fixed positive integer
The following property, which is part of the « folklore » (i.e., apparently unpublished as a specific result: see the introduction of this article in which it is stated as having « long been known ») has important consequences. For instance it rules out uniform distribution of the values of in the arithmetic progressions modulo for any integer .
For every fixed positive integer , the relation holds for almost all , meaning for all but values of as .
This is an elementary consequence of the fact that the sum of the reciprocals of the primes congruent to 1 modulo diverges, which itself is a corollary of the proof of Dirichlet's theorem on arithmetic progressions.
Generating functions
The Dirichlet series for may be written in terms of the Riemann zeta function as:
where the left-hand side converges for .
The Lambert series generating function is
which converges for .
Both of these are proved by elementary series manipulations and the formulae for .
Growth rate
In the words of Hardy & Wright, the order of is "always 'nearly '."
First
but as n goes to infinity, for all
These two formulae can be proved by using little more than the formulae for and the divisor sum function .
In fact, during the proof of the second formula, the inequality
true for , is proved.
We also have
Here is Euler's constant, , so and .
Proving this does not quite require the prime number theorem. Since goes to infinity, this formula shows that
In fact, more is true.
and
The second inequality was shown by Jean-Louis Nicolas. Ribenboim says "The method of proof is interesting, in that the inequality is shown first under the assumption that the Riemann hypothesis is true, secondly under the contrary assumption."
For the average order, we have
due to Arnold Walfisz, its proof exploiting estimates on exponential sums due to I. M. Vinogradov and N. M. Korobov.
By a combination of van der Corput's and Vinogradov's methods, H.-Q. Liu (On Euler's function.Proc. Roy. Soc. Edinburgh Sect. A 146 (2016), no. 4, 769–775)
improved the error term to
(this is currently the best known estimate of this type). The "Big " stands for a quantity that is bounded by a constant times the function of inside the parentheses (which is small compared to ).
This result can be used to prove that the probability of two randomly chosen numbers being relatively prime is .
Ratio of consecutive values
In 1950 Somayajulu proved
In 1954 Schinzel and Sierpiński strengthened this, proving that the set
is dense in the positive real numbers. They also proved that the set
is dense in the interval (0,1).
Totient numbers
A totient number is a value of Euler's totient function: that is, an for which there is at least one for which . The valency or multiplicity of a totient number is the number of solutions to this equation. A nontotient is a natural number which is not a totient number. Every odd integer exceeding 1 is trivially a nontotient. There are also infinitely many even nontotients, and indeed every positive integer has a multiple which is an even nontotient.
The number of totient numbers up to a given limit is
for a constant .
If counted accordingly to multiplicity, the number of totient numbers up to a given limit is
where the error term is of order at most for any positive .
It is known that the multiplicity of exceeds infinitely often for any .
Ford's theorem
proved that for every integer there is a totient number of multiplicity : that is, for which the equation has exactly solutions; this result had previously been conjectured by Wacław Sierpiński, and it had been obtained as a consequence of Schinzel's hypothesis H. Indeed, each multiplicity that occurs, does so infinitely often.
However, no number is known with multiplicity . Carmichael's totient function conjecture is the statement that there is no such .
Perfect totient numbers
A perfect totient number is an integer that is equal to the sum of its iterated totients. That is, we apply the totient function to a number n, apply it again to the resulting totient, and so on, until the number 1 is reached, and add together the resulting sequence of numbers; if the sum equals n, then n is a perfect totient number.
Applications
Cyclotomy
In the last section of the Disquisitiones Gauss proves that a regular -gon can be constructed with straightedge and compass if is a power of 2. If is a power of an odd prime number the formula for the totient says its totient can be a power of two only if is a first power and is a power of 2. The primes that are one more than a power of 2 are called Fermat primes, and only five are known: 3, 5, 17, 257, and 65537. Fermat and Gauss knew of these. Nobody has been able to prove whether there are any more.
Thus, a regular -gon has a straightedge-and-compass construction if n is a product of distinct Fermat primes and any power of 2. The first few such are
2, 3, 4, 5, 6, 8, 10, 12, 15, 16, 17, 20, 24, 30, 32, 34, 40,... .
Prime number theorem for arithmetic progressions
The RSA cryptosystem
Setting up an RSA system involves choosing large prime numbers and , computing and , and finding two numbers and such that . The numbers and (the "encryption key") are released to the public, and (the "decryption key") is kept private.
A message, represented by an integer , where , is encrypted by computing .
It is decrypted by computing . Euler's Theorem can be used to show that if , then .
The security of an RSA system would be compromised if the number could be efficiently factored or if could be efficiently computed without factoring .
Unsolved problems
Lehmer's conjecture
If is prime, then . In 1932 D. H. Lehmer asked if there are any composite numbers such that divides . None are known.
In 1933 he proved that if any such exists, it must be odd, square-free, and divisible by at least seven primes (i.e. ). In 1980 Cohen and Hagis proved that and that . Further, Hagis showed that if 3 divides then and .
Carmichael's conjecture
This states that there is no number with the property that for all other numbers , , . See Ford's theorem above.
As stated in the main article, if there is a single counterexample to this conjecture, there must be infinitely many counterexamples, and the smallest one has at least ten billion digits in base 10.
Riemann hypothesis
The Riemann hypothesis is true if and only if the inequality
is true for all where is Euler's constant and is the product of the first primes.
| Mathematics | Specific functions | null |
53453 | https://en.wikipedia.org/wiki/Fermat%27s%20little%20theorem | Fermat's little theorem | In number theory, Fermat's little theorem states that if is a prime number, then for any integer , the number is an integer multiple of . In the notation of modular arithmetic, this is expressed as
For example, if and , then , and is an integer multiple of .
If is not divisible by , that is, if is coprime to , then Fermat's little theorem is equivalent to the statement that is an integer multiple of , or in symbols:
For example, if and , then , and is a multiple of .
Fermat's little theorem is the basis for the Fermat primality test and is one of the fundamental results of elementary number theory. The theorem is named after Pierre de Fermat, who stated it in 1640. It is called the "little theorem" to distinguish it from Fermat's Last Theorem.
History
Pierre de Fermat first stated the theorem in a letter dated October 18, 1640, to his friend and confidant Frénicle de Bessy. His formulation is equivalent to the following:
If is a prime and is any integer not divisible by , then is divisible by .
Fermat's original statement was
This may be translated, with explanations and formulas added in brackets for easier understanding, as:
Every prime number [] divides necessarily one of the powers minus one of any [geometric] progression [] [that is, there exists such that divides ], and the exponent of this power [] divides the given prime minus one [divides ]. After one has found the first power [] that satisfies the question, all those whose exponents are multiples of the exponent of the first one satisfy similarly the question [that is, all multiples of the first have the same property].
Fermat did not consider the case where is a multiple of nor prove his assertion, only stating:
(And this proposition is generally true for all series [sic] and for all prime numbers; I would send you a demonstration of it, if I did not fear going on for too long.)
Euler provided the first published proof in 1736, in a paper titled "Theorematum Quorundam ad Numeros Primos Spectantium Demonstratio" (in English: "Demonstration of Certain Theorems Concerning Prime Numbers") in the Proceedings of the St. Petersburg Academy, but Leibniz had given virtually the same proof in an unpublished manuscript from sometime before 1683.
The term "Fermat's little theorem" was probably first used in print in 1913 in Zahlentheorie by Kurt Hensel:
(There is a fundamental theorem holding in every finite group, usually called Fermat's little theorem because Fermat was the first to have proved a very special part of it.)
An early use in English occurs in A.A. Albert's Modern Higher Algebra (1937), which refers to "the so-called 'little' Fermat theorem" on page 206.
Further history
Some mathematicians independently made the related hypothesis (sometimes incorrectly called the Chinese hypothesis) that if and only if is prime. Indeed, the "if" part is true, and it is a special case of Fermat's little theorem. However, the "only if" part is false: For example, , but 341 = 11 × 31 is a pseudoprime to base 2. See below.
Proofs
Several proofs of Fermat's little theorem are known. It is frequently proved as a corollary of Euler's theorem.
Generalizations
Euler's theorem is a generalization of Fermat's little theorem: For any modulus and any integer coprime to , one has
where denotes Euler's totient function (which counts the integers from 1 to that are coprime to ). Fermat's little theorem is indeed a special case, because if is a prime number, then .
A corollary of Euler's theorem is: For every positive integer , if the integer is coprime with , then
for any integers and .
This follows from Euler's theorem, since, if , then for some integer , and one has
If is prime, this is also a corollary of Fermat's little theorem. This is widely used in modular arithmetic, because this allows reducing modular exponentiation with large exponents to exponents smaller than .
Euler's theorem is used with not prime in public-key cryptography, specifically in the RSA cryptosystem, typically in the following way: if
retrieving from the values of , and is easy if one knows . In fact, the extended Euclidean algorithm allows computing the modular inverse of modulo , that is, the integer such that
It follows that
On the other hand, if is the product of two distinct prime numbers, then . In this case, finding from and is as difficult as computing (this has not been proven, but no algorithm is known for computing without knowing ). Knowing only , the computation of has essentially the same difficulty as the factorization of , since , and conversely, the factors and are the (integer) solutions of the equation .
The basic idea of RSA cryptosystem is thus: If a message is encrypted as , using public values of and , then, with the current knowledge, it cannot be decrypted without finding the (secret) factors and of .
Fermat's little theorem is also related to the Carmichael function and Carmichael's theorem, as well as to Lagrange's theorem in group theory.
Converse
The converse of Fermat's little theorem fails for Carmichael numbers. However, a slightly weaker variant of the converse is Lehmer's theorem:
If there exists an integer such that
and for all primes dividing one has
then is prime.
This theorem forms the basis for the Lucas primality test, an important primality test, and Pratt's primality certificate.
Pseudoprimes
If and are coprime numbers such that is divisible by , then need not be prime. If it is not, then is called a (Fermat) pseudoprime to base . The first pseudoprime to base 2 was found in 1820 by Pierre Frédéric Sarrus: 341 = 11 × 31.
A number that is a Fermat pseudoprime to base for every number coprime to is called a Carmichael number. Alternately, any number satisfying the equality
is either a prime or a Carmichael number.
Miller–Rabin primality test
The Miller–Rabin primality test uses the following extension of Fermat's little theorem:
If is an odd prime and with and odd > 0, then for every coprime to , either or there exists such that and .
This result may be deduced from Fermat's little theorem by the fact that, if is an odd prime, then the integers modulo form a finite field, in which 1 modulo has exactly two square roots, 1 and −1 modulo .
Note that holds trivially for , because the congruence relation is compatible with exponentiation. And holds trivially for since is odd, for the same reason. That is why one usually chooses a random in the interval .
The Miller–Rabin test uses this property in the following way: given an odd integer for which primality has to be tested, write with and odd > 0, and choose a random such that ; then compute ; if is not 1 nor −1, then square it repeatedly modulo until you get −1 or have squared times. If and −1 has not been obtained by squaring, then is a composite and is a witness for the compositeness of . Otherwise, is a strong probable prime to base a; that is, it may be prime or not. If is composite, the probability that the test declares it a strong probable prime anyway is at most , in which case is a strong pseudoprime, and is a strong liar. Therefore after non-conclusive random tests, the probability that is composite is at most 4−k, and may thus be made as low as desired by increasing .
In summary, the test either proves that a number is composite or asserts that it is prime with a probability of error that may be chosen as low as desired. The test is very simple to implement and computationally more efficient than all known deterministic tests. Therefore, it is generally used before starting a proof of primality.
| Mathematics | Modular arithmetic | null |
53455 | https://en.wikipedia.org/wiki/Minkowski%27s%20theorem | Minkowski's theorem | In mathematics, Minkowski's theorem is the statement that every convex set in which is symmetric with respect to the origin and which has volume greater than contains a non-zero integer point (meaning a point in that is not the origin). The theorem was proved by Hermann Minkowski in 1889 and became the foundation of the branch of number theory called the geometry of numbers. It can be extended from the integers to any lattice and to any symmetric convex set with volume greater than , where denotes the covolume of the lattice (the absolute value of the determinant of any of its bases).
Formulation
Suppose that is a lattice of determinant in the -dimensional real vector space and is a convex subset of that is symmetric with respect to the origin, meaning that if is in then is also in . Minkowski's theorem states that if the volume of is strictly greater than , then must contain at least one lattice point other than the origin. (Since the set is symmetric, it would then contain at least three lattice points: the origin 0 and a pair of points , where .)
Example
The simplest example of a lattice is the integer lattice of all points with integer coefficients; its determinant is 1. For , the theorem claims that a convex figure in the Euclidean plane symmetric about the origin and with area greater than 4 encloses at least one lattice point in addition to the origin. The area bound is sharp: if is the interior of the square with vertices then is symmetric and convex, and has area 4, but the only lattice point it contains is the origin. This example, showing that the bound of the theorem is sharp, generalizes to hypercubes in every dimension .
Proof
The following argument proves Minkowski's theorem for the specific case of
Proof of the case: Consider the map
Intuitively, this map cuts the plane into 2 by 2 squares, then stacks the squares on top of each other. Clearly has area less than or equal to 4, because this set lies within a 2 by 2 square. Assume for a contradiction that could be injective, which means the pieces of cut out by the squares stack up in a non-overlapping way. Because is locally area-preserving, this non-overlapping property would make it area-preserving for all of , so the area of would be the same as that of , which is greater than 4. That is not the case, so the assumption must be false: is not injective, meaning that there exist at least two distinct points in that are mapped by to the same point: .
Because of the way was defined, the only way that can equal is for
to equal for some integers and , not both zero.
That is, the coordinates of the two points differ by two even integers.
Since is symmetric about the origin, is also a point in . Since is convex, the line segment between and lies entirely in , and in particular the midpoint of that segment lies in . In other words,
is a point in . But this point is an integer point, and is not the origin since and are not both zero.
Therefore, contains a nonzero integer point.
Remarks:
The argument above proves the theorem that any set of volume contains two distinct points that differ by a lattice vector. This is a special case of Blichfeldt's theorem.
The argument above highlights that the term is the covolume of the lattice .
To obtain a proof for general lattices, it suffices to prove Minkowski's theorem only for ; this is because every full-rank lattice can be written as for some linear transformation , and the properties of being convex and symmetric about the origin are preserved by linear transformations, while the covolume of is and volume of a body scales by exactly under an application of .
Applications
Bounding the shortest vector
Minkowski's theorem gives an upper bound for the length of the shortest nonzero vector. This result has applications in lattice cryptography and number theory.
Theorem (Minkowski's bound on the shortest vector): Let be a lattice. Then there is a with . In particular, by the standard comparison between and norms, .
Remarks:
The constant in the bound can be improved, for instance by taking the open ball of radius as in the above argument. The optimal constant is known as the Hermite constant.
The bound given by the theorem can be very loose, as can be seen by considering the lattice generated by . But it cannot be further improved in the sense that there exists a global constant such that there exists an -dimensional lattice satisfying for all . Furthermore, such lattice can be self-dual.
Even though Minkowski's theorem guarantees a short lattice vector within a certain magnitude bound, finding this vector is in general a hard computational problem. Finding the vector within a factor guaranteed by Minkowski's bound is referred to as Minkowski's Vector Problem (MVP), and it is known that approximation SVP reduces to it using transference properties of the dual lattice. The computational problem is also sometimes referred to as HermiteSVP.
The LLL-basis reduction algorithm can be seen as a weak but efficiently algorithmic version of Minkowski's bound on the shortest vector. This is because a -LLL reduced basis for has the property that ; see these lecture notes of Micciancio for more on this. As explained in, proofs of bounds on the Hermite constant contain some of the key ideas in the LLL-reduction algorithm.
Applications to number theory
Primes that are sums of two squares
The difficult implication in Fermat's theorem on sums of two squares can be proven using Minkowski's bound on the shortest vector.
Theorem: Every prime with can be written as a sum of two squares.
Additionally, the lattice perspective gives a computationally efficient approach to Fermat's theorem on sums of squares:
First, recall that finding any nonzero vector with norm less than in , the lattice of the proof, gives a decomposition of as a sum of two squares. Such vectors can be found efficiently, for instance using LLL-algorithm. In particular, if is a -LLL reduced basis, then, by the property that , . Thus, by running the LLL-lattice basis reduction algorithm with , we obtain a decomposition of as a sum of squares. Note that because every vector in has norm squared a multiple of , the vector returned by the LLL-algorithm in this case is in fact a shortest vector.
Lagrange's four-square theorem
Minkowski's theorem is also useful to prove Lagrange's four-square theorem, which states that every natural number can be written as the sum of the squares of four natural numbers.
Dirichlet's theorem on simultaneous rational approximation
Minkowski's theorem can be used to prove Dirichlet's theorem on simultaneous rational approximation.
Algebraic number theory
Another application of Minkowski's theorem is the result that every class in the ideal class group of a number field contains an integral ideal of norm not exceeding a certain bound, depending on , called Minkowski's bound: the finiteness of the class number of an algebraic number field follows immediately.
Complexity theory
The complexity of finding the point guaranteed by Minkowski's theorem, or the closely related Blichfeldt's theorem, have been studied from the perspective of TFNP search problems. In particular, it is known that a computational analogue of Blichfeldt's theorem, a corollary of the proof of Minkowski's theorem, is PPP-complete. It is also known that the computational analogue of Minkowski's theorem is in the class PPP, and it was conjectured to be PPP complete.
| Mathematics | Other | null |
53457 | https://en.wikipedia.org/wiki/Quinoa | Quinoa | Quinoa (Chenopodium quinoa; , from Quechua or ) is a flowering plant in the amaranth family. It is a herbaceous annual plant grown as a crop primarily for its edible seeds; the seeds are rich in protein, dietary fiber, B vitamins and dietary minerals in amounts greater than in many grains. Quinoa is not a grass but rather a pseudocereal botanically related to spinach and amaranth (Amaranthus spp.), and originated in the Andean region of northwestern South America. It was first used to feed livestock 5,2007,000 years ago, and for human consumption 3,0004,000 years ago in the Lake Titicaca basin of Peru and Bolivia.
The plant thrives at high elevations and produces seeds that are rich in protein. Almost all production in the Andean region is done by small farms and associations. Its cultivation has spread to more than 70 countries, including Kenya, India, the United States, and European countries. As a result of increased consumption in North America, Europe, and Australasia, quinoa crop prices tripled between 2006 and 2014, entering a boom and bust cycle.
The quinoa monoculture that arose from increased production, combined with climate change effects in the native Andean region, created challenges for production and yield, and led to environmental degradation.
Description
Chenopodium quinoa is a dicotyledonous annual plant, usually about high. It has broad, generally powdery, hairy, lobed leaves, normally arranged alternately. The woody central stem is branched or unbranched depending on the variety and may be green, red or purple. The flowering panicles arise from the top of the plant or from leaf axils along the stem. Each panicle has a central axis from which a secondary axis emerges either with flowers (amaranthiform) or bearing a tertiary axis carrying the flowers (glomeruliform). These are small, incomplete, sessile flowers of the same colour as the sepals, and both pistillate and perfect forms occur. Pistillate flowers are generally located at the proximal end of the glomeruli and the perfect ones at the distal end of it. A perfect flower has five sepals, five anthers and a superior ovary, from which two to three stigmatic branches emerge.
The green hypogynous flowers have a simple perianth and are generally self-fertilizing, though cross-pollination occurs. In the natural environment, betalains serve to attract animals to generate a greater rate of pollination and ensure, or improve, seed dissemination. The fruits (seeds) are about in diameter and of various colors — from white to red or black, depending on the cultivar.
In regards to the "newly" developed salinity resistance of C.quinoa, some studies have concluded that accumulation of organic osmolytes plays a dual role for the species. They provide osmotic adjustment, in addition to protection against oxidative stress of the photosynthetic structures in developing leaves. Studies also suggested that reduction in stomatal density in reaction to salinity levels represents an essential instrument of defence to optimize water use efficiency under the given conditions to which it may be exposed.
Taxonomy
The species Chenopodium quinoa was first described by Carl Ludwig Willdenow (1765–1812), a German botanist who studied plants from South America, brought back by explorers Alexander von Humboldt and Aimé Bonpland.
Quinoa is an allotetraploid plant, containing two full sets of chromosomes from two different species which hybridised with each other at one time. According to a 1979 study, its presumed ancestor is either Chenopodium berlandieri, from North America, or the Andean species Ch. hircinum, although more recent studies, in 2011, even suggest Old World relatives. On the other hand, morphological features relate Ch.quinoa of the Andes and Ch. nuttalliae of Mexico. Some studies have suggested that both species may have been derived from the same wild type. A weedy quinoa, Ch.quinoa var. melanospermum, is known from South America, but no equivalent closely related to Ch.nutalliae has been reported from Mexico so far.
Studies regarding the genetic diversity of quinoa suggest that it may have passed through at least three bottleneck genetic events, with a possible fourth expected:
The first occurred when the species was created, as its two diploid ancestors underwent a hybridization followed by chromosome doubling, this new species was genetically isolated from its parent species, and thus lost a great deal of genetic diversity. These ancestors are still not known, but are not the higher altitude crop species Chenopodium pallidicaule (cañahua), a diploid.
A second bottleneck may have occurred when quinoa was domesticated from its unknown but possible wild tetraploid form. It might have been domesticated twice: once in the high Andes and a second time in the Chilean and Argentinean lowlands.
A third bottleneck can be considered "political", and has lasted more than 400 years, from the Spanish conquest of the new continent until the present time. During this phase quinoa has been replaced with maize, marginalized from production processes possibly due to its important medicinal, social and religious roles for the indigenous populations of South America, but also because it is very difficult to process (dehusk) compared with maize.
In the 21st century, a fourth bottleneck event may occur, as traditional farmers migrate from rural zones to urban centers, which exposes quinoa to the risk of further genetic erosion. Better breeding may also result in loss of genetic diversity, as breeders would be expected to reduce unwanted alleles to produce uniform cultivars, but cross-breeding between local landraces has and will likely produce high-diversity cultivars.
Etymology
The genus name Chenopodium is composed of two words coming from the Greek χήν,-νός, goose and πόδῖον, podion "little foot", or "goose foot", because of the resemblance of the leaves with the trace of a goose's foot.
The specific epithet quinoa is a borrowing from the Spanish quinua or quinoa, itself derived from Quechua kinuwa.
The Incas nicknamed quinoa chisiya mama, which in Quechua means "mother of all grains".
Distribution
Chenopodium quinoa is believed to have been domesticated in the Peruvian Andes from wild or weed populations of the same species. There are non-cultivated quinoa plants (Chenopodium quinoa var. melanospermum) that grow in the area it is cultivated; these may either be related to wild predecessors, or they could be descendants of cultivated plants.
Cultivation
Over the last 5,000years the biogeography of Ch. quinoa has changed greatly, mainly by human influence, convenience and preference. It has changed not only in the area of distribution, but also in regards to the climate this plant was originally adapted to, in contrast to the climates on which it is able to successfully grow in now. In a process started by a number of pre-Inca South American indigenous cultures, people in Chile have been adapting quinoa to salinity and other forms of stress over the last 3,000years. Quinoa is also cultivated, since an early date, near the coast of northern Chile, where it was grown by the Chinchorro culture. Ch. quinoa was brought to the lowlands of south-central Chile at an early date from the Andean highlands. Varieties in the lowlands of south-central Chile derive directly from ancestral cultivars which then evolved in parallel to those of the highlands. It has been suggested that the introduction of Ch. quinoa occurred before highland varieties with floury perisperm emerged. There are wide discrepancies in the suggested dates of introduction, one study suggests c. 1000 BC as the introduction date while another suggests 600–1100 AD. In colonial times the plant is known to have been cultivated as far south as the Chiloé Archipelago and the shores of Nahuel Huapi Lake. The cuisine of Chiloé included bread made of Quinoa until at least the mid-19th century.
In Chile it had almost disappeared by the early 1940s; as of 2015 the crop is mostly grown in three areas by only some 300 smallholder farmers. Each of these areas is different: indigenous small-scale growers near the border with Bolivia who grow many types of Bolivian forms, a few farmers in the central region who exclusively grow a white-seeded variety and generally market their crops through a well-known cooperative, and in the south by women in home gardens in Mapuche reserves.
When Amaranthaceae became abundant in Lake Pacucha, Peru, the lake was fresh, and the lack of Amaranthaceae taxa strongly indicates droughts which turned the lake into a saltmarsh. Based on the pollen associated with soil manipulation, this is an area of the Andes where domestication of C.quinoa became popular, although it was not the only one. It was domesticated in various geographical zones. With this, morphological adaptations began to happen until having five ecotypes today. Quinoa's genetic diversity illustrates that it was and is a vital crop.Andean agronomists and nutrition scientists began researching quinoa in the early twentieth century, and it became the subject of much interest among researchers involved in neglected and underutilized crop studies in the 1970s.
In 2004, the international community became increasingly interested in quinoa and it entered a boom and bust economic cycle that would last for over ten years. Between 2004 and 2011, quinoa became a more interesting commodity and global excitement for it increased. At this point, Bolivia and Peru were the only major producers of quinoa. In 2013, there was an extreme increase in imports of quinoa by the United States, Canada and various European countries. In 2016, growth began to slow. Imports were still increasing but at a slower rate and quinoa prices declined as other countries began producing it. By 2015, over 75 countries were producing quinoa, as opposed to only eight countries in the 1980s.
Particularly for the high variety of Chilean landraces, in addition to how the plant has adapted to different latitudes, this crop is now potentially cultivable almost anywhere in the world.
Climate requirements
The plant's growth is highly variable due to the number of different subspecies, varieties and landraces (domesticated plants or animals adapted to the environment in which they originated). However, it is generally undemanding and altitude-hardy; it is grown from coastal regions to over in the Andes near the equator, with most of the cultivars being grown between and . Depending on the variety, optimal growing conditions are in cool climates with temperatures that vary between during the night to near during the day. Some cultivars can withstand lower temperatures without damage. Light frosts normally do not affect the plants at any stage of development, except during flowering. Midsummer frosts during flowering, a frequent occurrence in the Andes, lead to sterilization of the pollen. Rainfall requirements are highly variable between the different cultivars, ranging from during the growing season. Growth is optimal with well-distributed rainfall during early growth and no rain during seed maturation and harvesting.
United States
Quinoa has been cultivated in the United States, primarily in the high elevation San Luis Valley of Colorado where it was introduced in 1983. In this high-altitude desert valley, maximum summer temperatures rarely exceed and night temperatures are about . In the 2010s, experimental production was attempted in the Palouse region of Eastern Washington, and farmers in Western Washington began producing the crop. The Washington State University Skagit River Valley research facility near Mount Vernon grew thousands of its own experimental varieties. The Puget Sound region's climate is similar to that of coastal Chile where the crop has been grown for centuries. Due to the short growing season, North American cultivation requires short-maturity varieties, typically of Bolivian origin. Quinoa is planted in Idaho where a variety developed and bred specifically for the high-altitude Snake River Plain is the largest planted variety in North America.
Europe
Several countries within Europe have successfully grown quinoa on a commercial scale. Southern England, Holland and Denmark all have significant production.
Sowing
Quinoa requires a significant amount of precipitation in order to germinate, therefore the traditional sowing date in Peru was between September and November. To increase the chance that more crops survive it would be advantageous to split up the sowing date among the plants. Traditionally quinoa was sowed by broadcast, in rows or grooves, or by broadcast and then making rows. Soil preparation should occur before sowing, and weeding should come soon after sowing the seeds.
Rotation is used in its Andean native range. Rotation is common with potato, cereals and legumes including Lupinus mutabilis. Traditionally, quinoa rotation happens in plots called aynoqas. These are made up of different sized plots in different zones, and each family unit would own plots in different areas. The aynoqas allowed for better crop yield, agricultural and ecological sustainability, and food security within communities.
Soil
Quinoa plants do best in sandy, well-drained soils with a low nutrient content, moderate salinity, and a soil pH of 6 to8.5. The seedbed must be well prepared and drained to avoid waterlogging.
Quinoa has gained attention for its adaptability to contrasting environments such as saline soils, nutrient-poor soils and drought stressed marginal agroecosystems.
Genetics
The genome of quinoa was sequenced in 2017. Through traditional selective breeding and, potentially, genetic engineering, the plant is being modified to have higher crop yield, improved tolerance to heat and biotic stress, and greater sweetness through saponin inhibition.
Harvesting
Traditionally, quinoa grain is harvested by hand, and only rarely by machine, because the extreme variability of the maturity period of most quinoa cultivars complicates mechanization. Harvest needs to be precisely timed to avoid high seed losses from shattering, and different panicles on the same plant mature at different times. The crop yield in the Andean region (often around 3 t/ha up to 5 t/ha) is comparable to wheat yields. In the United States, varieties have been selected for uniformity of maturity and are mechanically harvested using conventional small grain combines.
Processing
The plants are allowed to stand until the stalks and seeds have dried out and the grain has reached a moisture content below 10%.
Handling involves threshing the seedheads from the chaff and winnowing the seed to remove the husk. Before storage, the seeds need to be dried in order to avoid germination. This was traditionally done manually, which is labour-intensive.
Production
In 2020, world production of quinoa was 175,188 tonnes, led by Peru and Bolivia with 97% of the total when combined.
Price
Since the early 21st century when quinoa became more commonly consumed in North America, Europe, and Australasia where it was not typically grown, the crop value increased. Between 2006 and 2013, quinoa crop prices tripled. In 2011, the average price was US$3,115 per tonne with some varieties selling as high as $8,000 per tonne. This compares with wheat prices of about US$340 per tonne, making wheat about 10% of the value of quinoa. The resulting effect on traditional production regions in Peru and Bolivia also influenced new commercial quinoa production elsewhere in the world, such as the United States. By 2013, quinoa was being cultivated in some 70 countries. As a result of expanding production outside the Andean highlands native for quinoa, the price plummeted starting in early 2015 and remained low for years. From 2018 to 2019, quinoa production in Peru declined by 22%. Some refer to this as the "quinoa bust" because of the devastation the price fall caused for farmers and industry.
Effects of rising demand on growers
Rising quinoa prices over the period of 2006 to 2017 may have reduced the affordability of quinoa to traditional consumers. However, a 2016 study using Peru's Encuesta Nacional de Hogares found that rising quinoa prices during 2004–2013 led to net economic benefits for producers, and other commentary indicated similar conclusions, including for women specifically. It has also been suggested that as quinoa producers rise above subsistence-level income, they switch their own consumption to Western processed foods which are often less healthy than a traditional, quinoa-based diet, whether because quinoa is held to be worth too much to keep for oneself and one's family, or because processed foods have higher status despite their poorer nutritional value. Efforts are being made in some areas to distribute quinoa more widely and ensure that farming and poorer populations have access to it and have an understanding of its nutritional importance, including use in free school breakfasts and government provisions distributed to pregnant and nursing women in need.
In terms of wider social consequences, research on traditional producers in Bolivia has emphasised a complex picture. The degree to which individual producers benefit from the global quinoa boom depends on its mode of production, for example through producer associations and co-operatives such as the Asociación Nacional de Productores de Quinua (founded in the 1970s), contracting through vertically integrated private firms, or wage labor. State regulation and enforcement may promote a shift to cash-cropping among some farmers and a shift toward subsistence production among others, while enabling many urban refugees to return to working the land, outcomes with complex and varied social effects.
The growth of quinoa consumption outside of its indigenous region has raised concerns over food security of the original consumers, unsustainably intensive farming of the crop, expansion of farming into otherwise marginal agricultural lands with concurrent loss of the natural environment, threatening both the sustainability of producer agriculture and the biodiversity of quinoa. Studies have found that smallholder traditional farming of quinoa, specifically in the Andean region of Peru has significantly less of an environmental impact in carbon produced, than the modern industrial quinoa production.
World demand for quinoa is sometimes presented in the media particularly as being caused by rising veganism, but one academic has commented that despite the drawbacks of quinoa, meat production in most cases is still less sustainable than quinoa.
Monoculture and climate change impacts
Because of the increasing demand for quinoa, some fields in the Andean region of Bolivia have become quinoa monocultures. Particularly in the Uyuni salt flats, soil degradation has occurred due to mechanized production and decreased vegetation cover after clearing for quinoa fields. This degradation has led to poorer quinoa yields and lower environmental health in the region.
Signs of desertification of the landscape is amplified by the effects of climate change on quinoa fields and the salt flats. Drier and hotter weather negatively affects quinoa production, while also increasing pest populations attacking quinoa and reducing the nutrient quality of the soil.
Quinoa became a grain of growing interest partially due to its ability to withstand many different climate conditions. Its native Andean region is prone to dry and wet spells, and to cold and hot temperatures. Research shows that quinoa prefers warmer temperatures and alternating irrigation. The randomness of weather conditions due to climate change has hindered development of quinoa crops.
The quinoa boom and bust cycle led to a periodic increased demand for quinoa which originally resulted in increased production in its native area. However, when other countries recognized the economic benefit of producing quinoa, its cultivation in Europe and the United States increased. Some studies indicate that it may be more productive to grow quinoa in the United States, particularly in Washington State, and in China rather than in its native regions.
Chemistry
In their natural state, the seeds have a coating that contains bitter-tasting saponins, making them unpalatable. Most of the grain sold commercially has been processed to remove this coating. This bitterness has beneficial effects during cultivation, as it deters birds and, therefore, the plant requires minimal protection. The genetic control of bitterness involves quantitative inheritance. Although lowering the saponin content through selective breeding to produce sweeter, more palatable varieties is complicated by ≈10% cross-pollination, it is a major goal of quinoa breeding programs, which may include genetic engineering.
The toxicity category rating of the saponins in quinoa treats them as mild eye and respiratory irritants and as a low gastrointestinal irritant. In South America, these saponins have many uses, including as a detergent for clothing and washing, and as a folk medicine antiseptic for skin injuries.
Additionally, the leaves and stems of all species of the genus Chenopodium and related genera of the family Amaranthaceae, including quinoa, contain high levels of oxalic acid.
Uses
The increasing demand for quinoa is partially due to the attention it received as a food that may help alleviate food insecurity in some world regions. Quinoa is high in protein, which makes it a possible alternative to meat for vegetarians and vegans, and for people who are lactose intolerant. It also has high concentrations of dietary minerals. Quinoa does not contain gluten. Some of these qualities may have improved the market to economically privileged people in North America, possibly increasing the price of quinoa.
Quinoa is an important food for the Indigenous people of the Andean Altiplano, especially the Aymara and Quechua communities. Historically, it was consumed as a subsistence food, which was devalued by the Spanish when they colonized the region. The Spanish noticed that quinoa was consumed everyday and as a part of special ceremonies, so they decided it could grant power to people and threatened their conquest. Because of this they targeted it for extinction and significantly reduced the range where quinoa was grown.
For the Indigenous communities, growing quinoa represented food security and well-being, and it was involved in almost every meal of the day.
Nutrition
Raw, uncooked quinoa is 13% water, 64% carbohydrates, 14% protein, and 6% fat. Nutritional evaluations indicate that a serving of raw quinoa seeds has a food energy of and is a rich source (20% or higher of the Daily Value, DV) of protein, dietary fiber, several B vitamins, including 46%DV for folate, and for several dietary minerals such as magnesium (55%DV), manganese (95%DV), phosphorus (65%DV), and zinc (33%DV).
After boiling, which is the typical preparation for eating the seeds, many nutritional evaluations change. Although a serving of cooked quinoa increases to 72% water, most nutritional evaluations are reduced, such as, 21% carbohydrates, 4% protein, and 2% fat, and the food energy of cooked quinoa is reduced to . Although similarly reduced, cooked quinoa remains a rich source of the dietary minerals manganese (30%DV) and phosphorus (22%DV). However, cooked quinoa is reduced to being just a moderate source (10–19%DV) of dietary fiber and folate (11%), as well as of the dietary minerals iron (11%), magnesium (18%), and zinc (11%).
Quinoa is gluten-free. Because quinoa has a high concentration of protein and is a good source of many micronutrients, has versatility in preparation, and a potential for increased yields in controlled environments, it has been selected as an experimental crop in NASA's Controlled Ecological Life Support System for long-duration human occupied space flights.
In culture
United Nations recognition
The United Nations General Assembly declared 2013 as the "", in recognition of the ancestral practices of the Andean people, who have preserved it as a food for present and future generations, through knowledge and practices of living in harmony with nature. The objective was to draw the world's attention to the role that quinoa could play in providing food security, nutrition and poverty eradication in support of achieving Millennium Development Goals. Some academic commentary emphasized that quinoa production could have ecological and social drawbacks in its native regions, and that these problems needed to be tackled.
Kosher certification
Quinoa is used in the Jewish community as a substitute for the leavened grains that are forbidden during the Passover holiday. Several kosher certification organizations refuse to certify it as being kosher for Passover, citing reasons including its resemblance to prohibited grains or fear of cross-contamination of the product from nearby fields of prohibited grain or during packaging. However, in December 2013 the Orthodox Union, the world's largest kosher certification agency, announced it would begin certifying quinoa as kosher for Passover.
Gallery
| Biology and health sciences | Caryophyllales | null |
53460 | https://en.wikipedia.org/wiki/Statcoulomb | Statcoulomb | The statcoulomb (statC), franklin (Fr), or electrostatic unit of charge (esu) is the unit of measurement for electrical charge used in the centimetre–gram–second electrostatic units variant (CGS-ESU) and Gaussian systems of units. In terms of the Gaussian base units, it is
That is, it is defined so that the proportionality constant in Coulomb's law using CGS-ESU quantities is a dimensionless quantity equal to 1.
Definition and relation to CGS base units
Coulomb's law in the CGS-Gaussian system takes the form
where F is the force, q and q are the two electric charges, and r is the distance between the charges. This serves to define charge as a quantity in the Gaussian system.
The statcoulomb is defined such that if two electric charges of 1 statC each and have a separation of , the force of mutual electrical repulsion is 1 dyne. Substituting F = 1 dyn, q = q = 1 statC, and r = 1 cm, we get:
From this it is also evident that the quantity dimension of electric charge as defined in the CGS-ESU and Gaussian systems is .
Conversion between systems
Conversion of a quantity to the corresponding quantity of the International System of Quantities (ISQ) that underlies the International System of Units (SI) by using the defining equations of each system.
The SI uses the coulomb (C) as its unit of electric charge. The conversion factor between corresponding quantities with the units coulomb and statcoulomb depends on which quantity is to be converted. The most common cases are:
For electric charge:
For electric flux ():
For electric flux density ():
The symbol "≘" ('corresponds to') is used instead of "=" because the two sides cannot be equated.
| Physical sciences | Charge | Basics and measurement |
53464 | https://en.wikipedia.org/wiki/Cinnabar | Cinnabar | Cinnabar (; ), or cinnabarite (), also known as mercurblende is the bright scarlet to brick-red form of mercury(II) sulfide (HgS). It is the most common source ore for refining elemental mercury and is the historic source for the brilliant red or scarlet pigment termed vermilion and associated red mercury pigments.
Cinnabar generally occurs as a vein-filling mineral associated with volcanic activity and alkaline hot springs. The mineral resembles quartz in symmetry and it exhibits birefringence. Cinnabar has a mean refractive index near 3.2, a hardness between 2.0 and 2.5, and a specific gravity of approximately 8.1. The color and properties derive from a structure that is a hexagonal crystalline lattice belonging to the trigonal crystal system, crystals that sometimes exhibit twinning.
Cinnabar has been used for its color since antiquity in the Near East, including as a rouge-type cosmetic, in the New World since the Olmec culture, and in China since as early as the Yangshao culture, where it was used in coloring stoneware. In Roman times, cinnabar was highly valued as paint for walls, especially interiors (since it darkened when used outdoors due to exposure to sunlight).
Associated modern precautions for the use and handling of cinnabar arise from the toxicity of the mercury component, which was recognized as early as ancient Rome.
Etymology
The name comes from Greek (), a Greek word most likely applied by Theophrastus to several distinct substances. In Latin, it was sometimes known as minium, meaning also "red cinnamon", though both of these terms now refer specifically to lead tetroxide.
Properties and structure
Properties
Cinnabar is generally found in a massive, granular, or earthy form and is bright scarlet to brick-red in color, though it occasionally occurs in crystals with a nonmetallic adamantine luster. It resembles quartz in its symmetry. It exhibits birefringence, and it has the second-highest refractive index of any mineral. Its mean refractive index is 3.08 (sodium light wavelengths), versus the indices for diamond and the non-mineral gallium(III) arsenide (GaAs), which are 2.42 and 3.93, respectively. The hardness of cinnabar is 2.0–2.5 on the Mohs scale, and its specific gravity 8.1.
Structure
Structurally, cinnabar belongs to the trigonal crystal system. It occurs as thick tabular or slender prismatic crystals or as granular to massive incrustations. Crystal twinning occurs as simple contact twins.
Mercury(II) sulfide, HgS, adopts the cinnabar structure described, and one additional structure, i.e. it is dimorphous. Cinnabar is the more stable form, and is a structure akin to that of HgO: each Hg center has two short Hg−S bonds (each 2.36 Å), and four longer contacts (with 3.10, 3.10, 3.30 and 3.30 Å separations). In addition, HgS is found in a black, non-cinnabar polymorph (metacinnabar) that has the zincblende structure.
Occurrence
Cinnabar generally occurs as a vein-filling mineral associated with volcanic activity and alkaline hot springs. Cinnabar is deposited by epithermal ascending aqueous solutions (those near the surface and not too hot) far removed from their igneous source. It is associated with native mercury, stibnite, realgar, pyrite, marcasite, opal, quartz, chalcedony, dolomite, calcite, and barite.
Cinnabar is found in essentially all mineral extraction localities that yield mercury, notably Almadén (Spain). This mine was exploited from Roman times until 1991, being for centuries the most important cinnabar deposit in the world. Good cinnabar crystals have also been found there. Cinnabar deposits appear in Giza (Egypt); Puerto Princesa (Philippines); Red Devil, Alaska; Murfreesboro, Arkansas; New Almaden Mine in San Jose, California; New Idria, California, the Hastings Mine and St. John's Mine both in Vallejo, California; Terlingua, Texas (United States); Idrija (Slovenia); near Obermoschel in the Palatinate; the La Ripa and Levigliani mines at the foot of the Apuan Alps and in Mount Amiata (Tuscany, Italy); Avala (Serbia); Huancavelica (Peru); the province of Guizhou in China and Western ghats in India where fine crystals have been obtained. It has been found in Dominica near its sulfur springs at the southern end of the island along the west coast.
Cinnabar is still being deposited, such as from the hot waters of Sulphur Bank Mine in California and Steamboat Springs, Nevada (United States).
Mining and extraction of mercury
As the most common source of mercury in nature, cinnabar has been mined for thousands of years, even as far back as the Neolithic Age. During the Roman Empire it was mined both as a pigment, and for its mercury content.
To produce liquid mercury (quicksilver), crushed cinnabar ore is roasted in rotary furnaces. Pure mercury separates from sulfur in this process and easily evaporates. A condensing column is used to collect the liquid metal, which is most often shipped in iron flasks.
Toxicity
Associated modern precautions for use and handling of cinnabar arise from the toxicity of the mercury component, which was recognized as early as in ancient Rome. Because of its mercury content, cinnabar can be toxic to human beings. Overexposure to mercury, mercury poisoning (mercurialism), was seen as an occupational disease to the ancient Romans. Though people in ancient South America often used cinnabar for art, or processed it into refined mercury (as a means to gild silver and gold to objects), the toxic properties of mercury were well known. It was dangerous to those who mined and processed cinnabar; it caused shaking, loss of sense, and death. Data suggests that mercury was retorted from cinnabar and the workers were exposed to the toxic mercury fumes. "Mining in the Spanish cinnabar mines of Almadén, southwest of Madrid, was regarded as being akin to a death sentence due to the shortened life expectancy of the miners, who were slaves or convicts."
Decorative use
Cinnabar has been used for its color since antiquity in the Near East, including as a rouge-type cosmetic, in the New World since the Olmec culture, and in China for writing on oracle bones as early as the Zhou dynasty. Late in the Song dynasty it was used in coloring lacquerware.
Cinnabar's use as a color in the New World, since the Olmec culture, is exemplified by its use in royal burial chambers during the peak of Maya civilization, most dramatically in the 7th-century tomb of the Red Queen in Palenque, where the remains of a noble woman and objects belonging to her in her sarcophagus were completely covered with bright red powder made from cinnabar.
The most popularly known use of cinnabar is in Chinese carved lacquerware, a technique that apparently originated in the Song dynasty. The danger of mercury poisoning may be reduced in ancient lacquerware by entraining the powdered pigment in lacquer, but could still pose an environmental hazard if the pieces were accidentally destroyed. In the modern jewellery industry, the toxic pigment is replaced by a resin-based polymer that approximates the appearance of pigmented lacquer.
Two female mummies dated AD 1399 to 1475 found in Cerro Esmeralda in Chile in 1976 had clothes colored with cinnabar.
Other forms
Hepatic cinnabar, or paragite, is an impure brownish variety from the mines of Idrija in the Carniola region of Slovenia, in which the cinnabar is mixed with bituminous and earthy matter.
Hypercinnabar crystallizes at high temperature in the hexagonal crystal system.
Metacinnabar is a black-colored form of mercury(II) sulfide, which crystallizes in the cubic crystal system.
Synthetic cinnabar is produced by treatment of mercury(II) salts with hydrogen sulfide to precipitate black, synthetic metacinnabar, which is then heated in water. This conversion is promoted by the presence of sodium sulfide.
| Physical sciences | Minerals | Earth science |
53469 | https://en.wikipedia.org/wiki/Cinnamon | Cinnamon | Cinnamon is a spice obtained from the inner bark of several tree species from the genus Cinnamomum. Cinnamon is used mainly as an aromatic condiment and flavouring additive in a wide variety of cuisines, sweet and savoury dishes, breakfast cereals, snack foods, bagels, teas, hot chocolate and traditional foods. The aroma and flavour of cinnamon derive from its essential oil and principal component, cinnamaldehyde, as well as numerous other constituents including eugenol.
Cinnamon is the name for several species of trees and the commercial spice products that some of them produce. All are members of the genus Cinnamomum in the family Lauraceae. Only a few Cinnamomum species are grown commercially for spice. Cinnamomum verum (alternatively C. zeylanicum), known as "Ceylon cinnamon" after its origins in Sri Lanka (formerly Ceylon), is considered to be "true cinnamon", but most cinnamon in international commerce is derived from four other species, usually and more correctly referred to as "cassia": C. burmanni (Indonesian cinnamon or Padang cassia), C. cassia (Chinese cinnamon or Chinese cassia), C. loureiroi (Saigon cinnamon or Vietnamese cassia), and the less common C. citriodorum (Malabar cinnamon).
In 2021, world production of cinnamon was 226,753 tonnes, led by China with 43% of the total.
Etymology
The English word "cinnamon", attested in English since the 15th century, deriving from the Ancient Greek (, later κίνναμον : ), via Latin and medieval French intermediate forms. The Greek was borrowed from a Phoenician word, which was similar to the related Hebrew word ().
The name "cassia", first recorded in late Old English from Latin, ultimately derives from the Hebrew word , a form of the verb , "to strip off bark".
Early Modern English also used the names canel and canella, similar to the current names of cinnamon in several other European languages, which are derived from the Latin word , a diminutive of , "tube", from the way the bark curls up as it dries.
History
Cinnamon has been known from remote antiquity. It was imported to Egypt as early as 2000 BC, but those who reported that it had come from China had confused it with Cinnamomum cassia, a related species. Cinnamon was so highly prized among ancient nations that it was regarded as a gift fit for monarchs and even for a deity; an inscription records the gift of cinnamon and cassia to the temple of Apollo at Miletus. Its source was kept a trade secret in the Mediterranean world for centuries by those in the spice trade, in order to protect their monopoly as suppliers.
Cinnamomum verum, which translates from Latin as "true cinnamon", is native to India, Sri Lanka, Bangladesh and Myanmar. Cinnamomum cassia (cassia) is native to China. Related species, all harvested and sold in the modern era as cinnamon, are native to Vietnam ("Saigon cinnamon"), Indonesia and other southeast Asian countries with warm climates.
In Ancient Egypt, cinnamon was used to embalm mummies. From the Ptolemaic Kingdom onward, Ancient Egyptian recipes for kyphi, an aromatic used for burning, included cinnamon and cassia. The gifts of Hellenistic rulers to temples sometimes included cassia and cinnamon.
The first Greek reference to is found in a poem by Sappho in the 7th century BC. According to Herodotus, both cinnamon and cassia grew in Arabia, together with incense, myrrh and , and were guarded by winged serpents. Herodotus, Aristotle and other authors named Arabia as the source of cinnamon; they recounted that giant "cinnamon birds" collected the cinnamon sticks from an unknown land where the cinnamon trees grew and used them to construct their nests.
Pliny the Elder wrote that cinnamon was brought around the Arabian Peninsula on "rafts without rudders or sails or oars", taking advantage of the winter trade winds. He also mentioned cassia as a flavouring agent for wine, and that the tales of cinnamon being collected from the nests of cinnamon birds was a traders' fiction made up to charge more. However, the story remained current in Byzantium as late as 1310.
According to Pliny the Elder, a Roman pound () of cassia, cinnamon (), cost up to 1,500 , the wage of fifty months' labour. Diocletian's Edict on Maximum Prices from 301 AD gives a price of 125 for a pound of cassia, while an agricultural labourer earned 25 per day. Cinnamon was too expensive to be commonly used on funeral pyres in Rome, but the Emperor Nero is said to have burned a year's worth of the city's supply at the funeral for his wife Poppaea Sabina in AD 65.
Middle Ages
Through the Middle Ages, the source of cinnamon remained a mystery to the Western world. From reading Latin writers who quoted Herodotus, Europeans had learned that cinnamon came up the Red Sea to the trading ports of Egypt, but where it came from was less than clear. When the Sieur de Joinville accompanied his king, Louis IX of France to Egypt on the Seventh Crusade in 1248, he reported—and believed—what he had been told: that cinnamon was fished up in nets at the source of the Nile out at the edge of the world (i.e., Ethiopia). Marco Polo avoided precision on the topic.
The first mention that the spice grew in the area of India was in Maimonides's Mishneh Torah, about 1180. The first mention that the spice grew specifically in Sri Lanka was in Zakariya al-Qazwini's ("Monument of Places and History of God's Bondsmen") about 1270. This was followed shortly thereafter by John of Montecorvino in a letter of about 1292.
Indonesian rafts transported cinnamon directly from the Moluccas to East Africa (see also Rhapta), where local traders then carried it north to Alexandria in Egypt. Venetian traders from Italy held a monopoly on the spice trade in Europe, distributing cinnamon from Alexandria. The disruption of this trade by the rise of other Mediterranean powers, such as the Mamluk sultans and the Ottoman Empire, was one of many factors that led Europeans to search more widely for other routes to Asia.
Early modern period
During the 1500s, Ferdinand Magellan was searching for spices on behalf of Spain; in the Philippines, he found , which was closely related to C. zeylanicum, the cinnamon found in Sri Lanka. This cinnamon eventually competed with Sri Lankan cinnamon, which was controlled by the Portuguese.
In 1638, Dutch traders established a trading post in Sri Lanka, took control of the manufactories by 1640, and expelled the remaining Portuguese by 1658. "The shores of the island are full of it," a Dutch captain reported, "and it is the best in all the Orient. When one is downwind of the island, one can still smell cinnamon eight leagues out to sea." The Dutch East India Company continued to overhaul the methods of harvesting in the wild and eventually began to cultivate its own trees.
In 1767, Lord Brown of the British East India Company established the Anjarakkandy Cinnamon Estate near Anjarakkandy in the Kannur district of Kerala, India. It later became Asia's largest cinnamon estate. The British took control of Ceylon from the Dutch in 1796.
Cultivation
Cinnamon is an evergreen tree characterized by oval-shaped leaves, thick bark and a berry fruit. When harvesting the spice, the bark and leaves are the primary parts of the plant used. However, in Japan, the more pungent roots are harvested in order to produce nikki (ニッキ) which is a product distinct from cinammon (シナモン shinamon). Cinnamon is cultivated by growing the tree for two years, then coppicing it, i.e., cutting the stems at ground level. The following year, about a dozen new shoots form from the roots, replacing those that were cut. A number of pests such as Colletotrichum gloeosporioides, Diplodia species and Phytophthora cinnamomi (stripe canker) can affect the growing plants.
The stems must be processed immediately after harvesting while the inner bark is still wet. The cut stems are processed by scraping off the outer bark, then beating the branch evenly with a hammer to loosen the inner bark, which is then pried off in long rolls. Only of the inner bark is used; the outer, woody portion is discarded, leaving metre-long cinnamon strips that curl into rolls ("quills") on drying. The processed bark dries completely in four to six hours, provided it is in a well-ventilated and relatively warm environment. Once dry, the bark is cut into lengths for sale.
A less than ideal drying environment encourages the proliferation of pests in the bark, which may then require treatment by fumigation with sulphur dioxide. In 2011, the European Union approved the use of sulphur dioxide at a concentration of up to for the treatment of C. verum bark harvested in Sri Lanka.
Species
A number of species are often sold as cinnamon:
Cinnamomum cassia (cassia or Chinese cinnamon, the most common commercial type in the USA)
C. burmanni (Korintje, Padang cassia, or Indonesian cinnamon)
C. loureiroi (Saigon cinnamon, Vietnamese cassia, or Vietnamese cinnamon)
C. verum (Sri Lanka cinnamon, Ceylon cinnamon or Cinnamomum zeylanicum)
C. citriodorum (Malabar cinnamon)
Cassia induces a strong, spicy flavour and is often used in baking, especially associated with cinnamon rolls, as it handles baking conditions well. Among cassia, Chinese cinnamon is generally medium to light reddish-brown in colour, hard and woody in texture, and thicker ( thick), as all of the layers of bark are used. Ceylon cinnamon, using only the thin inner bark, has a lighter brown colour and a finer, less dense, and more crumbly texture. It is subtle and more aromatic in flavour than cassia and it loses much of its flavour during cooking.
The barks of the species are easily distinguished when whole, both in macroscopic and microscopic characteristics. Ceylon cinnamon sticks (quills) have many thin layers and can easily be made into powder using a coffee or spice grinder, whereas cassia sticks are much harder. Indonesian cinnamon is often sold in neat quills made up of one thick layer, capable of damaging a spice or coffee grinder. Saigon cinnamon (C. loureiroi) and Chinese cinnamon (C. cassia) are always sold as broken pieces of thick bark, as the bark is not supple enough to be rolled into quills.
The powdered bark is harder to distinguish, but if it is treated with tincture of iodine (a test for starch), little effect is visible with pure Ceylon cinnamon; however, when Chinese cinnamon is present, a deep-blue tint is produced.
Grading
The Sri Lankan grading system divides the cinnamon quills into four groups:
Alba, less than in diameter
Continental, less than in diameter
Mexican, less than in diameter
Hamburg, less than in diameter
These groups are further divided into specific grades. For example, Mexican is divided into M00000 special, M000000 and M0000, depending on quill diameter and number of quills per kilogram. Any pieces of bark less than long are categorized as quillings. Featherings are the inner bark of twigs and twisted shoots. Chips are trimmings of quills, outer and inner bark that cannot be separated, or the bark of small twigs.
Production
In 2021, four countries accounted for 98% of the world's cinnamon production, a total of 226,753 tonnes: China, Indonesia, Vietnam, and Sri Lanka.
Counterfeit
True cinnamon from C. verum bark can be mixed with cassia (C. cassia) as counterfeit and falsely marketed as authentic cinnamon. In one analysis, authentic Ceylon cinnamon bark contained 12-143 mg/kg of coumarin a phenolic typically low in content in true cinnamon but market samples contained coumarin with levels as high as 3462 mg/kg, indicating probable contamination with cassia in the counterfeit cinnamon. ConsumerLab.com found the same problem in a 2020 analysis; "a supplement that contained the highest amount of coumarin was labeled as Ceylon cinnamon".
Food uses
Cinnamon bark is used as a spice. It is principally employed in cookery as a condiment and flavouring material. It is used in the preparation of chocolate, especially in Mexico. Cinnamon is often used in savoury dishes of chicken and lamb. In the United States and Europe, cinnamon and sugar are often used to flavour cereals, bread-based dishes such as toast, and fruits, especially apples; a cinnamon and sugar mixture (cinnamon sugar) is sold separately for such purposes. It is also used in Portuguese and Turkish cuisine for both sweet and savoury dishes. Cinnamon can also be used in pickling, and in Christmas drinks such as eggnog. Cinnamon powder has long been an important spice in enhancing the flavour of Persian cuisine, used in a variety of thick soups, drinks and sweets.
Cinnamon is a common ingredient in Jewish cuisine across various communities. In Sephardic cooking, it is incorporated into vegetable stews and desserts such as tishpishti and travados, both of which are soaked in honey. In Ashkenazi cuisine, cinnamon features in dishes like honey cakes, and kugels. It is also one of "four sibling spices" (rempah empat beradik) essential in Malay cuisine along with clove, star anise and cardamom.
Nutrient composition
Ground cinnamon is 11% water, 81% carbohydrates (including 53% dietary fiber), 4% protein and 1% fat.
Characteristics
Texture
Ceylon cinnamon may be crushed into small pieces by hand while Indonesian cinnamon requires a powerful blender.
Flavour, aroma and taste
The flavour of cinnamon is due to the aromatic essential oils that makes up 0.5 to 1% of its composition.
Cinnamon bark can be macerated, then extracted in 80% ethanol, to a tincture.
Cinnamon essential oil can be prepared by roughly pounding the bark, macerating it in sea water, and then quickly distilling the whole. It is of a golden-yellow colour, with the characteristic odour of cinnamon and a very hot aromatic taste.
Cinnamon oil nanoemulsion can be made with polysorbate 80, cinnamon essential oil, and water, by ultrasonic emulsification.
Cinnamon oil macroemulsion can be made with a dispersing emulsifying homogenizer.
The pungent taste and scent come from cinnamaldehyde, about 90% of the essential oil from cinnamon bark. Cinnamaldehyde decomposes, in high humidity and high temperatures, to styrene, and, by reaction with oxygen as it ages, it darkens in colour and forms resinous compounds.
Cinnamon constituents include some 80 aromatic compounds, including eugenol, found in the oil from leaves or bark of cinnamon trees.
Alcohol flavorant
Cinnamon is used as a flavoring in cinnamon liqueur, such as cinnamon-flavored whiskey in the United States, and , a cinnamon brandy in Greece.
Health-related research
Cinnamon has a long history of use in traditional medicine as a digestive aid. However, contemporary studies are unable to find evidence of any significant medicinal or therapeutic effect.
Reviews of clinical trials reported lowering of fasting plasma glucose and inconsistent effects on hemoglobin A1C (HbA1c, an indicator of chronically elevated plasma glucose). Four of the reviews reported a decrease in fasting plasma glucose, only two reported lower HbA1c, and one reported no change to either measure. The Cochrane review noted that trial durations were limited to 4 to 16 weeks, and that no trials reported on changes to quality of life, morbidity or mortality rate. The Cochrane authors' conclusion was: "There is insufficient evidence to support the use of cinnamon for type 1 or type 2 diabetes mellitus." Citing the Cochrane review, the U.S. National Center for Complementary and Integrative Health stated: "Studies done in people don't support using cinnamon for any health condition." However, the results of the studies are difficult to interpret because it is often unclear what type of cinnamon and what part of the plant were used.
A meta-analysis of cinnamon supplementation trials with lipid measurements reported lower total cholesterol and triglycerides, but no significant changes in LDL-cholesterol or HDL-cholesterol. Another reported no change to body weight or insulin resistance.
Toxicity
A systematic review of adverse events as a result of cinnamon use reported gastrointestinal disorders and allergic reactions as the most frequently reported side effects.
In 2008, the European Food Safety Authority considered the toxicity of coumarin, a component of cinnamon, and confirmed a maximum recommended tolerable daily intake (TDI) of 0.1 mg of coumarin per kg of body weight. Coumarin is known to cause liver and kidney damage in high concentrations and metabolic effect in humans with CYP2A6 polymorphism. Based on this assessment, the European Union set a guideline for maximum coumarin content in foodstuffs of 50 mg per kg of dough in seasonal foods, and 15 mg per kg in everyday baked foods. The maximum recommended TDI of 0.1 mg of coumarin per kg of body weight equates to 5 mg of coumarin (or 5.6 g C. verum with 0.9 mg coumarin per gram) for a body weight of 50 kg. C as shown in the table below:
Due to the variable amount of coumarin in C. cassia, usually well over 1.0 mg of coumarin per g of cinnamon and sometimes up to 12 times that, C. cassia has a low safe-intake-level upper limit to adhere to the above TDI. In contrast, C. verum has only trace amounts of coumarin.
In March 2024, the US Food and Drug Administration recommended a voluntary recall on 6 brands of cinnamon due to contamination with lead, after an investigation stemming from 500 reports of child lead poisoning across the US. The FDA determined that cinnamon was adulterated with lead chromate.
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53479 | https://en.wikipedia.org/wiki/2060%20Chiron | 2060 Chiron | 2060 Chiron is a ringed small Solar System body in the outer Solar System, orbiting the Sun between Saturn and Uranus. Discovered in 1977 by Charles Kowal, it was the first-identified member of a new class of objects now known as centaurs—bodies orbiting between the asteroid belt and the Kuiper belt. Chiron is named after the centaur Chiron in Greek mythology.
Although it was initially called an asteroid and classified only as a minor planet with the designation "2060 Chiron", in 1989 it was found to exhibit behavior typical of a comet. Today it is classified as both a minor planet and a comet, and is accordingly also known by the cometary designation 95P/Chiron. More recently, a series of occultation events through the 2010s and early 2020s revealed that Chiron hosts rings, making it one of four minor planets and the only known comet to host a ring system.
History
Discovery
Chiron was discovered on 1 November 1977 by Charles Kowal from images taken on 18 October at Palomar Observatory. It was given the temporary designation of . It was found near aphelion and at the time of discovery it was the most distant known minor planet. Chiron was even claimed as the tenth planet by the press. Chiron was later found on several precovery images, going back to 1895, which allowed its orbit to be accurately determined. It had been at perihelion in 1945 but was not discovered then because there were few searches being made at that time, and these were not sensitive to slow-moving objects. The Lowell Observatory's survey for distant planets would not have gone down faint enough in the 1930s and did not cover the right region of the sky in the 1940s. The April 1895 precovery image was one month after the March 1895 perihelion.
Naming
This minor planet was named after Chiron, a half-human, half-horse centaur from Greek mythology. Son of the Titan Cronus and the nymph Philyra, Chiron was the wisest and most just of all centaurs, serving as an instructor of the Greek heroes. The official was published by the Minor Planet Center on 1 April 1978 (). It was suggested that the names of other centaurs be reserved for objects of the same type.
Chiron, along with most major and minor planetary bodies, is not generally given a symbol in astronomy. A symbol was devised for it by Al H. Morrison and is mostly used among astrologers: it resembles a key as well as an OK monogram for Object Kowal.
Orbit
Chiron's orbit was found to be highly eccentric (0.37), with perihelion just inside the orbit of Saturn and aphelion just outside the perihelion of Uranus (it does not reach the average distance of Uranus, however). According to the program Solex, Chiron's closest approach to Saturn in modern times was around May 720, when it came within million km () of the planet. During this passage Saturn's gravity caused Chiron's semi-major axis to decrease from AU to 13.7 AU. Chiron's orbit does not intersect Uranus' orbit.
Chiron attracted considerable interest because it was the first object discovered in such an orbit, well outside the asteroid belt. Chiron is classified as a centaur, the first of a class of objects orbiting between the outer planets. Chiron is a Saturn–Uranus object because its perihelion lies in Saturn's zone of control and its aphelion lies in that of Uranus. Centaurs are not in stable orbits and will be removed by gravitational perturbation by the giant planets over a period of millions of years, moving to different orbits or leaving the Solar System altogether. Chiron likely comes from the Kuiper belt and will probably become a short-period comet in about a million years. Chiron came to perihelion (closest point to the Sun) in 1996 and aphelion in May 2021.
Physical characteristics
Spectral type
The visible and near-infrared spectrum of Chiron is neutral, and is similar to that of C-type asteroids and the nucleus of Halley's Comet. The near-infrared spectrum of Chiron shows absence of water ice.
Rotation period
Four rotational light curves of Chiron were taken from photometric observations between 1989 and 1997. Lightcurve analysis gave a concurring, well-defined rotational period of 5.918 hours with a small brightness variation of 0.05 to 0.09 magnitude, which indicates that the body has a rather spheroidal shape ().
Diameter
The assumed size of an object depends on its absolute magnitude (H) and the albedo (the amount of light it reflects). In 1984 Lebofsky estimated Chiron to be about 180 km in diameter. Estimates in the 1990s were closer to 150 km in diameter. Occultation data from 1993 suggests a diameter of about 180 km. Combined data from the Spitzer Space Telescope in 2007 and the Herschel Space Observatory in 2011 suggests that Chiron is in diameter. Therefore, Chiron may be as large as 10199 Chariklo. The diameter of Chiron is difficult to estimate in part because the true absolute magnitude of its nucleus is uncertain due to its highly variable cometary activity.
Cometary behavior
In February 1988, at 12 AU from the Sun, Chiron brightened by 75 percent. This is behavior typical of comets but not asteroids. Further observations in April 1989 showed that Chiron had developed a cometary coma, A tail was detected in 1993. Chiron differs from other comets in that water is not a major component of its coma, because it is too far from the Sun for water to sublimate. In 1995 carbon monoxide was detected in Chiron in very small amounts, and the derived CO production rate was calculated to be sufficient to account for the observed coma. Cyanide was also detected in the spectrum of Chiron in 1991. At the time of its discovery, Chiron was close to aphelion, whereas the observations showing a coma were done closer to perihelion, perhaps explaining why no cometary behavior had been seen earlier. The fact that Chiron is still active probably means it has not been in its current orbit very long.
Chiron is officially designated as both a comet—95P/Chiron—and a minor planet, an indication of the sometimes fuzzy dividing line between the two classes of object. The term proto-comet has also been used. Being about 220 km in diameter, it is unusually large for a comet nucleus. Chiron was the first member of a new family of Chiron-type comets (CTC) with (TJupiter > 3; a > aJupiter). Other CTCs include: 39P/Oterma, 165P/LINEAR, 166P/NEAT, and 167P/CINEOS. There are also non-centaur asteroids that are simultaneously classified as comets, such as 4015 Wilson–Harrington, 7968 Elst–Pizarro, and 118401 LINEAR. Michael Brown lists it as possibly a dwarf planet with a measured diameter of , which may be near the lower limit for an icy object to have been a dwarf planet at some point in its history.
Since the discovery of Chiron, other centaurs have been discovered, and nearly all are currently classified as minor planets, but are being observed for possible cometary behavior. 60558 Echeclus has displayed a cometary coma and now also has the cometary designation 174P/Echeclus. After passing perihelion in early 2008, 52872 Okyrhoe significantly brightened.
Rings
Chiron has rings, similar to the better-established rings of 10199 Chariklo. Based on unexpected occultation events observed in stellar-occultation data obtained on 7 November 1993, 9 March 1994, and 29 November 2011, which were initially interpreted as resulting from jets associated with Chiron's comet-like activity, Chiron's rings were proposed to be in radius and sharply defined. The rings' changing appearance at different viewing angles can largely explain the long-term variation in Chiron's brightness and hence estimates of Chiron's albedo and size. Moreover, it can, by assuming that the water ice is in Chiron's rings, explain the changing intensity of the infrared water-ice absorption bands in Chiron's spectrum, including their disappearance in 2001 (when the rings were edge-on). Also, the geometric albedo of Chiron's rings as determined by spectroscopy is consistent with that used to explain Chiron's long-term brightness variations.
Further evidence of the rings was provided by two independent observations of occultations on 28 November 2018 and 15 December 2022, which suggests that their structure is constantly evolving. In the 2018 event Chiron's rings were observed to have less material than in 2011, but seemed to be developing a partial third ring; by the 2022 event there was more material than either of the previous observations, and the third ring had become fully developed. J.L. Ortiz speculated that the extra material in the 2022 event could be from an outburst observed in 2021, which left more material in orbit and thus bolstered the generation of the third ring–this is also expected to be cyclical, maintaining the rings.
The preferred pole of Chiron's rings is, in ecliptic coordinates, λ = , β = . The rings' width, separation, and optical depths were observed to be nearly identical to those of Chariklo's rings until the 2018 observation, indicating that the same type of structure had been responsible for both. Moreover, both their rings are within their respective Roche limits, though Chiron's newly developed third ring may be outside of it depending on its density.
Exploration
The Chiron Orbiter Mission was a mission proposed for NASA's New Frontiers program or Flagship program. It was published in May 2010 and proposed an orbiter mission to Chiron. Its launch date could have varied from as early as 2023 to as late as 2025, depending on budget and propulsion type.
There was another mission proposed, part of the Discovery Program known as Centaurus; if approved, it would have launched between 2026 and 2029 and made a flyby of 2060 Chiron and one other Centaur sometime in the 2030s.
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| Physical sciences | Solar System | Astronomy |
53484 | https://en.wikipedia.org/wiki/Narwhal | Narwhal | The narwhal (Monodon monoceros) is a species of toothed whale native to the Arctic. It is the only member of the genus Monodon and one of two living representatives of the family Monodontidae. The narwhal is a stocky cetacean with a relatively blunt snout, a large melon, and a shallow ridge in place of a dorsal fin. Males of this species have a large () long tusk, which is a protruding left canine thought to function as a weapon, a tool for feeding, in attracting mates or sensing water salinity. Specially adapted slow-twitch muscles, along with the jointed neck vertebrae and shallow dorsal ridge allow for easy movement through the Arctic environment, where the narwhal spends extended periods at great depths. The narwhal's geographic range overlaps with that of the similarly built and closely related beluga whale, and the animals are known to interbreed.
Narwhals inhabits the Arctic waters of Canada, Greenland and Russia. Every year, they migrate to ice-free summering grounds, usually in shallow waters, and often return to the same sites in subsequent years. Their diet mainly consists of polar and Arctic cod, Greenland halibut, cuttlefish, shrimp, and armhook squid. Diving to depths of up to , the narwhal is among the deepest-diving cetaceans. The animals typically travel in groups of three to eight, with aggregations of up to 1,000 occurring in the summer months. Narwhals mate among the offshore pack ice from March to May, and the young are born between July and August of the following year. When communicating amongst themselves, narwhals use a variety of clicks, whistles and knocks.
There are an estimated 170,000 living narwhals, and the species is listed as being of least concern by the International Union for Conservation of Nature (IUCN). The population is threatened by the effects of climate change, such as reduction in ice cover and human activities such as pollution and hunting. Narwhals have been hunted for thousands of years by Inuit in northern Canada and Greenland for meat and ivory, and regulated subsistence hunting continues to this day.
Taxonomy
The narwhal was scientifically described by Carl Linnaeus in his 1758 publication Systema Naturae. The word "narwhal" comes from the Old Norse , meaning 'corpse-whale', which possibly refers to the animal's grey, mottled skin and its habit of remaining motionless when at the water's surface, a behaviour known as "logging" that usually happens in the summer. The scientific name, Monodon monoceros, is derived from Ancient Greek, meaning 'single-tooth single-horn'.
The narwhal is most closely related to the beluga whale (Delphinapterus leucas). Together, these two species comprise the only extant members of the family Monodontidae, sometimes referred to as the "white whales". Monodontids are distinguished by their pronounced melons (acoustic sensory organs), short snouts and the absence of a true dorsal fin.
Although the narwhal and beluga are classified as separate genera, there is some evidence of interbreeding between the two. Most prominent are the remains of a whale, described by marine zoologists as unlike any known species, which were found in West Greenland around 1990. It had features midway between a narwhal and a beluga, indicating that the remains belonged to a hybrid between the two species (a 'narluga'); this was confirmed by a 2019 DNA analysis. Whether the hybrid itself could breed remains unknown.
Evolution
Results of a genetic study reveal that porpoises and monodontids are closely related, forming a separate clade which diverged from other dolphins about 11 million years ago (mya). A 2018 molecular analysis of monodontid fossils indicates that they separated from Phocoenidae (porpoises) around 10.82 to 20.12 mya, and they are considered to be sister taxa. A later phylogenetic study conducted in 2020 suggested that the narwhal split from the beluga whale around 4.98 mya, based on data from mitochondrial DNA.
The fossil species Casatia thermophila of early Pliocene central Italy was described as a possible narwhal ancestor when it was discovered in 2019. Bohaskaia, Denebola and Haborodelphis are other extinct genera known from the Pliocene of the United States. Fossil evidence shows that prehistoric monodontids lived in tropical waters. They may have migrated to Arctic and subarctic waters in response to changes in the marine food chain.
The following phylogenetic tree is based on a 2019 study of the family Monodontidae.
Description
The narwhal has a robust body with a short, blunt snout, small upcurved flippers, and convex to concave tail flukes. Adults measure in length and weigh . Male narwhals attain sexual maturity at 12 to 20 years of age, reaching a length of . Females reach sexual maturity at a younger age, between 8 and 9 years old, when they are about long. On average, males are about longer and more than 75% heavier than females.
The colouration of the narwhal consists of a mottled pattern, with blackish-brown markings over a white background. At birth, the skin is light grey, and when sexually mature, white patches grow on the navel and genital slit; such whitening occurs throughout life, resulting in aged narwhals being almost purely white. Unlike most whales, the narwhal has a shallow dorsal ridge, rather than a dorsal fin, possibly an evolutionary adaptation to make swimming under ice easier or to facilitate rolling. The neck vertebrae are also jointed, instead of being fused as in most whales, which allows for a greater range of neck flexibility. These characteristics are shared by the beluga whale. Furthermore, male and female narwhals have differently shaped tail flukes; the former are bent inward, while the latter are swept back on the front margins. This is thought to be an adaptation for reducing drag caused by the tusk.
The skeletal muscles of narwhals are highly adapted for prolonged periods of deep-sea foraging. During such activities, oxygen is reserved in the muscles, which are typically slow-twitch, enabling greater endurance and manouverability. Narwhals also have a comparatively high amount of myoglobin in their body, which helps to facilitate deeper dives. It has a dense layer of blubber, around thick. This fat accounts for a third of the body mass and helps insulate from cold ocean temperatures.
Tusk and dentition
The most conspicuous trait of male narwhals is a long, spiralled tusk, which is a canine tooth that projects from the left side of the upper jaw. Both sexes have a pair of tusks embedded in the upper jaw, which in males erupt from the lip somewhere between two and three years of age. The tusk grows throughout the animal's life, reaching lengths of . It is hollow and weighs up to . Some males may grow two tusks, occurring when the right canine also protrudes through the lip. Females rarely grow tusks: when they do, the tusks are typically smaller than those of males, with less noticeable spirals.
Current scientific consensus indicates that narwhal tusks are secondary sexual characteristics which indicate social status. Further functions of the narwhal tusk are debated: while some biologists suggest that narwhals use their tusks in fights, others argue that they may be of use in feeding. The tusk is also a highly innervated sensory organ with millions of nerve endings, allowing the narwhal to sense temperature variability in its surroundings. These nerves may also be able to detect changes in particle concentration and water pressure. According to Dr. Martin Nweeia, male narwhals may rid themselves of encrustations on their tusks by rubbing them together, as opposed to posturing displays of aggressive male-to-male rivalry. Drone footage from August 2016 in Tremblay Sound, Nunavut, revealed that narwhals used their tusks to tap and stun small Arctic cod, making them easier to catch for feeding. Females, who usually do not have tusks, live longer than males, hence the tusk cannot be essential to the animal's survival. It is generally accepted that the primary function of the narwhal tusk is associated with sexual selection.
Alongside its tusk, the narwhal has a single pair of small vestigial teeth that reside in open tooth sockets in the upper jaw. These teeth, which differ in form and composition, encircle the exposed tooth sockets laterally, posteriorly, and ventrally. Vestigial teeth in male narwhals are commonly shed in the palate. The varied morphology and anatomy of small teeth indicate a path of evolutionary obsolescence.
Distribution
The narwhal is found in the Atlantic and Russian areas of the Arctic Ocean. Individuals are commonly recorded in the Canadian Arctic Archipelago, such as in the northern part of Hudson Bay, in Hudson Strait, in Baffin Bay, off the east coast of Greenland and in a strip running east from the northern end of Greenland to eastern Russia (170° east). Land in this strip includes Svalbard, Franz Joseph Land and Severnaya Zemlya. The northernmost sightings of narwhals occurred north of Franz Joseph Land, at about 85° north. There are an estimated 12,500 narwhals in the northern Hudson Bay, whereas around 140,000 reside in Baffin Bay.
Migration
Narwhals exhibit seasonal migration, with a high fidelity of return to preferred ice-free summering grounds, usually in shallow waters. In summer months, they move closer to the coast, often in pods of 10–100 individuals. In the winter, they move to deeper waters offshore, under thick pack ice, surfacing in narrow fissures or in wider fractures known as leads. As spring comes, these leads open up into channels and the narwhals return to the coastal bays. Narwhals in Baffin Bay typically travel to northern Canada and Greenland between June and September. After this period, they travel about south to the Davis Strait, and stay there until April. During winter, narwhals from Canada and West Greenland regularly visit the pack ice of the Davis Strait and Baffin Bay along the continental slope which contains less than 5% open water and hosts a high density of Greenland halibut.
Behaviour and ecology
Narwhals normally congregate in groups of three to eight individuals. Groups may be "nurseries" with only females and young, or can contain only juveniles or adult males ("bulls"); mixed groups can occur at any time of year. In the summer, several groups come together, forming larger aggregations which can contain 500 to over 1,000 individuals. Male narwhals have been observed rubbing each other's tusks, a behaviour known as "tusking".
When in their wintering waters, narwhals make some of the deepest dives recorded for cetaceans, diving to at least over 15 times per day, with many dives reaching . The greatest dive depth recorded is . Dives last up to 25 minutes, and vary in depth depending on the season and local variation between environments. For example, in the Baffin Bay wintering grounds, narwhals tend to dive deep within the steep coasts, typically south of Baffin Bay. This suggests differences in habitat structure, prey availability, or genetic adaptations between subpopulations. In the northern wintering grounds, narwhals do not dive as deep as the southern population, in spite of greater water depths in these areas. This is mainly attributed to prey being concentrated nearer to the surface, which causes narwhals to alter their foraging strategies.
Diet
Narwhals have a restricted and specialised diet. Due to the lack of well-developed dentition, narwhals are believed to feed by swimming close to prey and sucking them into the mouth. A study of the stomach contents of 73 narwhals found Arctic cod (Boreogadus saida) to be the most commonly consumed prey, followed by Greenland halibut (Reinhardtius hippoglossoides). Large quantities of Boreo-Atlantic armhook squid (Gonatus fabricii) were also discovered. Male specimens had a higher likelihood of showing two additional prey species within their stomach contents: polar cod (Arctogadus glacialis) and redfish (Sebastes marinus), both of which are found at depths of more than . The study also concluded that the size of prey did not differ between genders or age groups. Other items found within narwhal stomach contents include wolffish, capelin, skate eggs and sometimes rocks.
Narwhal diet varies between seasons. In winter, narwhals feed on demersal prey, mostly flatfish, under dense pack ice. During the summer, they eat mostly Arctic cod and Greenland halibut, with other fish such as polar cod making up the remainder of their diet. Narwhals consume more food in the winter months than they do in summer.
Breeding
Most female narwhals reproduce by the time they are six to eight years old. Courtship and mating behaviour for the species has been recorded from March to May, when they live among offshore pack ice, and is thought to involve a dominant male mating with several partners. The average gestation period lasts 15 months, and births appear to be most frequent between July and August. Female narwhals have a birth interval of around 2–3 years. As with most marine mammals, only a single calf is born, averaging in length with white or light grey pigmentation. Summer population surveys along different coastal inlets of Baffin Island found that calf numbers varied from 0.05% of 35,000 in Admiralty Inlet, to 5% of 10,000 total in Eclipse Sound. These findings suggest that higher calf counts may reflect calving and nursery habitats in favourable inlets.
Newborn calves begin their lives with a thin layer of blubber. The blubber thickens as they nurse their mother's milk, which is rich in fat; calves are dependent on milk for about 20 months. This long lactation period gives calves time to learn the skills they will need to survive as they mature. Narwhals are among the few animals that undergo menopause and live for decades after they have finished breeding. Females in this phase may continue to protect calves in the pod. A 2024 study concluded that five species of toothed whale evolved menopause to acquire higher overall longevity, although their reproductive periods did not change. To explain this, scientists hypothesised that calves of these species require the assistance of (post-)menopausal females for an enhanced chance at survival, as they are extremely difficult for a single female to successfully rear.
Communication
Like most toothed whales, narwhals use sound to navigate and hunt for food. They primarily vocalise through clicks, whistles and knocks, created by air movement between chambers near the blowhole. The frequency of these sounds ranges from 0.3 to 125 hertz, while those used for echolocation typically fall between 19 and 48 hertz. Sounds are reflected off the sloping front of the skull and focused by the animal's melon: a mass of fat which can be controlled through surrounding musculature. Echolocation clicks are used for detecting prey and locating barriers at short distances. Whistles and throbs are most commonly used to communicate with other pod members. Calls recorded from the same pod are more similar than calls from different pods, suggesting the possibility of group- or individual-specific calls. Narwhals sometimes adjust the duration and pitch of their pulsed calls to maximise sound propagation in varying acoustic environments. Other sounds produced by narwhals include trumpeting and "squeaking-door sounds". The narwhal vocal repertoire is similar to that of the beluga whale. However, the frequency ranges, durations, and repetition rates of narwhal clicks differ from those of belugas.
Longevity and mortality factors
Age determination techniques using the number of periosteum layers in the lower jaw reveal that narwhals live an average of 50 years, though techniques using amino acid dating from the lens of the eyes suggest that female narwhals can reach 115±10 years and male narwhals can live to 84±9 years.
Death by suffocation often occurs when narwhals fail to migrate before the Arctic freezes over in late autumn. This is known as "sea-ice entrapment". Narwhals drown if open water is no longer accessible and ice is too thick for them to break through. Breathing holes in ice may be up to apart, which limits the use of foraging grounds. These holes must be at least wide to allow an adult whale to breathe. Narwhals also die of starvation from entrapment events.
In 19141915, around 1,000 narwhal carcasses were discovered after entrapment events, most occurring in areas such as Disko Bay in West Greenland. Several cases of sea entrapment were recorded in 2008–2010, during the Arctic winter, including in some places where such events had never been recorded before. This suggests later departure dates from summering grounds. Wind and currents move sea ice from adjacent locations to Greenland, leading to fluctuations in concentration. Due to their tendency of returning to the same areas, changes in weather and ice conditions are not always associated with narwhal movement toward open water. It is currently unclear to what extent sea ice changes pose a danger to narwhals.
Narwhals are preyed upon by polar bears and orcas. In some instances, the former have been recorded waiting at breathing holes for young narwhals, while the latter were observed surrounding and killing entire narwhal pods. To escape predators such as orcas, narwhals may use prolonged submersion to hide under ice floes rather than relying on speed.
Researchers found bacteria of the Brucella genus in the bloodstreams of numerous narwhals throughout the course of a 19-year study. They were also recorded with whale lice species such as Cyamus monodontis and Cyamus nodosus. Other pathogens that affect narwhals include Toxoplasma gondii, morbillivirus, and papillomavirus. In 2018, a female narwhal was recorded with an alphaherpesvirus in her system.
Conservation
The narwhal is listed as a species of least concern by the IUCN Red List. As of 2017, the global population is estimated to be 123,000 mature individuals out of a total of 170,000. There were about 12,000 narwhals in Northern Hudson Bay in 2011, and around 49,000 near Somerset Island in 2013. There are approximately 35,000 in Admiralty Inlet, 10,000 in Eclipse Sound, 17,000 in Eastern Baffin Bay, and 12,000 in Jones Sound. Population numbers in Smith Sound, Inglefield Bredning and Melville Bay are 16,000, 8,000 and 3,000, respectively. There are roughly 800 narwhals in the waters off Svalbard.
In the 1972 Marine Mammal Protection Act, the United States banned imports of products made from narwhal parts. They are listed on Appendix II of the Convention on International Trade in Endangered Species of Wild Fauna and Flora (CITES) and Convention on the Conservation of Migratory Species of Wild Animals (CMS). These committees restrict international trading of live animals and their body parts, as well as implementing sustainable action plans. The species is classified as special concern under the Committee on the Status of Endangered Wildlife in Canada (COSEWIC), which aims to classify the risk levels of species in the country.
Threats
Narwhals are hunted for their skin, meat, teeth, tusks and carved vertebrae, which are commercially traded. About 1,000 narwhals are killed per year: 600 in Canada and 400 in Greenland. Canadian catches were steady at this level in the 1970s, dropped to 300–400 per year in the late 1980s and 1990s and have risen again since 1999. Greenland caught more, 700–900 per year, in the 1980s and 1990s.
In Canada and Greenland, Narwhal tusks are sold both carved and uncarved. Per hunted narwhal, an average of one to two vertebrae and teeth are sold. In Greenland, the skin () is sold commercially to fish factories, and in Canada to other communities. Based on an analysis of 2007 narwhal hunts in Hudson Bay, a 2013 paper estimated that gross revenue per narwhal was (US$). Hunts receive subsidies, but they continue mainly to support tradition, rather than for profit. Economic analysis noted that whale watching may be an alternate source of revenue.
As narwhals grow, bioaccumulation of heavy metals takes place within their bodies. It is thought that pollution in the ocean is the primary cause of bioaccumulation in marine mammals; this may lead to health problems for the narwhal population. When bioaccumulating, numerous metals appear in the blubber, liver, kidney and musculature. A study found that the blubber was nearly devoid of these metals, whereas the liver and kidneys had a dense concentration of them. Relative to the liver, the kidney has a greater concentration of zinc and cadmium, while lead, copper and mercury were not nearly as abundant. Individuals of different weight and sex showed differences in the concentration of metals in their organs.
Narwhals are one of the Arctic marine mammals most vulnerable to climate change due to sea ice decline, especially in their northern wintering grounds such as the Baffin Bay and Davis Strait regions. Satellite data collected from these areas shows the amount of sea ice has been markedly reduced from what it was previously. It is thought that narwhals' foraging ranges reflect patterns they acquired early in life, which improves their capacity to obtain the food supplies they need for the winter. This strategy focuses on strong site fidelity rather than individual-level responses to local prey distribution, resulting in focal foraging areas during the winter. As such, despite changing conditions, narwhals will continue to return to the same areas during migration.
Reduction in sea ice has possibly led to an increased exposure to predation. In 2002, hunters in Siorapaluk experienced an increase in the number of caught narwhals, but this increase did not seem to be linked to enhanced endeavour, implying that climate change may be making the narwhal more vulnerable to hunting. Scientists recommend assessing population numbers, assigning sustainable quotas, and ensuring local acceptance of sustainable development. Seismic surveys associated with oil exploration disrupt the narwhal's normal migration patterns. These disturbed migrations may also be associated with increased sea ice entrapment.
Relationship with humans
Narwhals have coexisted alongside circumpolar peoples for millennia. Their long, distinctive tusks were often held with fascination throughout human history. These tusks were prized for their supposed healing powers, and were worn on staffs and thrones. Depictions of narwhals in paintings such as The Lady and the Unicorn have found a prevalent place in human arts.
Inuit
Narwhals have been hunted by Inuit to the same extent as other sea mammals, such as seals and whales. Almost all parts of the narwhalthe meat, skin, blubber and organsare consumed. , the raw skin and attached blubber, is considered a delicacy. As a custom, one or two vertebrae per animal are used for tools and art. The skin is an important source of vitamin C, which is otherwise difficult to obtain in the Arctic Circle. In some places in Greenland, such as Qaanaaq, traditional hunting methods are used and whales are harpooned from handmade kayaks. In other parts of Greenland and Northern Canada, high-speed boats and hunting rifles are used.
In Inuit legend, the narwhal's tusk was created when a woman with harpoon rope tied around her waist was dragged into the ocean after the harpoon had stuck into a large narwhal. She was then transformed into a narwhal; her hair, which she was wearing in a twisted knot, became the spiralling narwhal tusk.
Tusk trade
In Europe, narwhal tusks were highly sought-after for centuries. This stems from a medieval belief that narwhal tusks were the horns of the legendary unicorn. Considered to have magical properties, narwhal tusks were used to counter poisoning, and all sorts of diseases such as measles and rubella. The rise of modern science towards the end of the 17th century led to a decreased belief in magic and alchemy. After the unicorn notion was scientifically refuted, narwhal tusks were rarely employed for magical purposes.
According to some theories, Vikings and Greenland Norse began trade of narwhal tusks, which, via European channels, would later reach markets in the Middle East and East Asia. The idea that Norsemen hunted narwhals was once held, but was never confirmed and is now considered improbable. Narwhal tusks were given as state gifts to kings and queens throughout medieval Europe, with the price of narwhal tusks said to have been a couple of hundred times greater than their weight in gold during the 18th and 19th centuries. Ivan the Terrible had a jewellery-covered narwhal tusk on his deathbed, while Elizabeth I received a narwhal tusk allegedly valued at £10,000 pounds sterling from the privateer Martin Frobisher. Both items were staples in cabinets of curiosities.
| Biology and health sciences | Toothed whale | Animals |
53539 | https://en.wikipedia.org/wiki/Soyuz%20programme | Soyuz programme | The Soyuz programme ( , ; , meaning "Union") is a human spaceflight programme initiated by the Soviet Union in the early 1960s. The Soyuz spacecraft was originally part of a Moon landing project intended to put a Soviet cosmonaut on the Moon. It was the third Soviet human spaceflight programme after the Vostok (1961–1963) and Voskhod (1964–1965) programmes.
The programme consists of the Soyuz capsule and the Soyuz rocket and is now the responsibility of the Russian Roscosmos. After the retirement of the Space Shuttle in 2011, Soyuz was the only way for humans to get to the International Space Station (ISS) until 30 May 2020, when Crew Dragon flew to the ISS for the first time with astronauts.
Soyuz rocket
The launch vehicles used in the Soyuz expendable launch system are manufactured at the Progress State Research and Production Rocket Space Center (TsSKB-Progress) in Samara, Russia. As well as being used in the Soyuz programme as the launcher for the crewed Soyuz spacecraft, Soyuz launch vehicles are now also used to launch robotic Progress supply spacecraft to the International Space Station and commercial launches marketed and operated by TsSKB-Progress and the Starsem company. Currently Soyuz vehicles are launched from the Baikonur Cosmodrome in Kazakhstan and the Plesetsk Cosmodrome in northwest Russia and, since 2011, Soyuz launch vehicles are also being launched from the Guiana Space Centre in French Guiana. The Spaceport's new Soyuz launch site has been handling Soyuz launches since 21 October 2011, the date of the first launch. As of December 2019, 19 Guiana Soyuz launches had been made from French Guiana Space Centre, all successful.
The Soyuz rocket family is one of the most dependable and widely utilized launch vehicles in the history of space travel. It has been in operation for nearly six decades, having been developed by the Soviet Union and presently run by Russia. The Soyuz rockets have played an important role in both crewed and uncrewed space missions, launching people to the International Space Station (ISS) and delivering satellites and scientific payloads.
Soyuz spacecraft
The basic Soyuz spacecraft design was the basis for many projects, many of which were never developed. Its earliest form was intended to travel to the Moon without employing a huge booster like the Saturn V or the Soviet N-1 by repeatedly docking with upper stages that had been put in orbit using the same rocket as the Soyuz. This and the initial civilian designs were done under the Soviet Chief Designer Sergei Pavlovich Korolev, who did not live to see the craft take flight. Several military derivatives took precedence in the Soviet design process, though they never came to pass.
A Soyuz spacecraft consists of three parts (from front to back):
a spheroid orbital module
a small aerodynamic reentry module
a cylindrical service module with solar panels attached
There have been many variants of the Soyuz spacecraft, including:
Sever early crewed spacecraft proposal to replace Vostok (1959)
L1-1960 crewed circumlunar spacecraft proposal (1960); evolved into the Soyuz-A design
L4-1960 crewed lunar orbiter proposal (1960)
L1-1962 crewed lunar flyby spacecraft proposal (1962); early design led to Soyuz
OS-1962 space station proposal (1962)
Soyuz-A 7K-9K-11K circumlunar complex proposal (1963)
Soyuz 7K crewed spacecraft concept; cancelled in 1964 in favor of the LK-1
Soyuz 9K proposed orbital tug; cancelled in 1964 when the Soyuz 7K and Soyuz P were cancelled
Soyuz 11K proposed fuel tanker; cancelled in 1964 when the Soyuz 7K and Soyuz P were cancelled
L3-1963 crewed lunar lander proposal (1963)
L4-1963 crewed lunar orbiter proposal; modified 7K (1963)
Soyuz 7K-OK (1967–1970)
Soyuz 7K-L1 Zond (1967–1970)
Soyuz 7K-L3 LOK (1971–1972)
Soyuz 7K-OKS (1971); also known as 7KT-OK
Soyuz 7K-T or "ferry" (1973–1981)
Soyuz 7K-T-AF (1973); 7K-T modified for space station flight with Orion 2 space telescope
Soyuz 7K-T/A9 (1974–1978); 7K-T modified for flights to military Almaz space stations
Soyuz 7K-TM (1974–1976)
7K-MF6 (1976); 7K-TM modified for space station flight with MKF-6 camera
Soyuz-T (1976–1986)
Zarya planned 'Super Soyuz' replacement for Soyuz and Progress (1985)
Alpha Lifeboat rescue spacecraft based on Zarya (1995); cancelled in favor of a modified Soyuz TM
Big Soyuz enlarged version of Soyuz reentry vehicle (2008)
Soyuz-TM (1986–2003)
Soyuz TMA (2003–2012)
Soyuz-ACTS (2006)
Soyuz TMA-M (2010–2016)
Soyuz MS (since 2016)
Military Soyuz (P, PPK, R, 7K-VI Zvezda, and OIS)
Soyuz P crewed satellite interceptor proposal (1962); cancelled in 1964 in favor of the Istrebitel Sputnikov program
Soyuz R command-reconnaissance spacecraft proposal (1962); cancelled in 1966 and replaced by Almaz
Soyuz 7K-TK transport spacecraft proposal for delivering cosmonauts to Soyuz R military stations (1966); cancelled in 1970 in favor of the TKS spacecraft
Soyuz PPK revised version of Soyuz P (1964)
Soyuz 7K-VI Zvezda space station proposal (1964)
Soyuz-VI crewed combat spacecraft proposal; cancelled in 1965
Soyuz OIS (1967)
Soyuz OB-VI space station proposal (1967)
Soyuz 7K-S military transport proposal (1974)
Soyuz 7K-ST concept for Soyuz T and TM (1974)
Derivatives
The Zond spacecraft was designed to take a crew around the Moon, but never achieved the required degree of safety or political need. Zond 5 did circle the Moon in September 1968, with two tortoises and other life forms, and returned safely to Earth although in an atmospheric entry which probably would have killed human travelers.
The Progress series of robotic cargo ships for the Salyut, Mir, and ISS use the engine section, orbital module, automatic navigation, docking mechanism, and overall layout of the Soyuz spacecraft, but are incapable of reentry.
While not a direct derivative, the Chinese Shenzhou spacecraft follows the basic template originally pioneered by Soyuz.
Soyuz crewed flights
Soviet human spaceflight missions started in 1961 and ended in 1991 with the dissolution of the Soviet Union.
The Russian human spaceflight missions program started in 1991 and continues to this day. Soyuz crewed missions were the only spacecraft visiting the International Space Station, starting from when the Space Shuttle program ended in 2011, until the launch of Crew Dragon Demo-2 on 30 May 2020. The International Space Station always has at least one Soyuz spacecraft docked at all times for use as an escape craft.
Soyuz uncrewed flights
Kosmos 133 - launch failure
Kosmos 140 - reentry damage
Kosmos 186
Kosmos 188
Kosmos 212
Kosmos 213
Kosmos 238
Soyuz 2 - failed to dock
Kosmos 379
Kosmos 396
Kosmos 434
Kosmos 496
Kosmos 573
Kosmos 613
Kosmos 638
Kosmos 656
Kosmos 670
Kosmos 672
Kosmos 772 - partial fail
Soyuz 20
Kosmos 869
Kosmos 1001
Kosmos 1074
Soyuz 34
Soyuz T-1
Soyuz TM-1
Soyuz MS-14
Soyuz MS-23
Gallery
| Technology | Programs and launch sites | null |
53551 | https://en.wikipedia.org/wiki/Charadriiformes | Charadriiformes | Charadriiformes (, from Charadrius, the type genus of family Charadriidae) is a diverse order of small to medium-large birds. It includes about 390 species and has members in all parts of the world. Most charadriiform birds live near water and eat invertebrates or other small animals; however, some are pelagic (seabirds), others frequent deserts, and a few are found in dense forest. Members of this group can also collectively be referred to as shorebirds.
Taxonomy, systematics and evolution
The order was formerly divided into three suborders:
The waders (or "Charadrii"): typical shorebirds, most of which feed by probing in the mud or picking items off the surface in both coastal and freshwater environments.
The gulls and their allies (or "Lari"): these are generally larger species which take fish from the sea. Several gulls and skuas will also take food items from beaches, or rob smaller species, and some have become adapted to inland environments.
The auks (or "Alcae") are coastal species which nest on sea cliffs and "fly" underwater to catch fish.
The Sibley-Ahlquist taxonomy lumps all the Charadriiformes together with other seabirds and birds of prey into a greatly enlarged order Ciconiiformes. However, the resolution of the DNA-DNA hybridization technique used by Sibley & Ahlquist was not sufficient to properly resolve the relationships in this group, and indeed it appears as if the Charadriiformes constitute a single large and very distinctive lineage of modern birds of their own.
The auks, usually considered distinct because of their peculiar morphology, are more likely related to gulls, the "distinctness" being a result of adaptation for diving.
Families
The order Charadriiformes contains 3 suborders, 19 families and 391 species.
Suborder Charadrii
Family Burhinidae – stone-curlews, thick-knees (10 species)
Family Pluvianellidae – Magellanic plover
Family Chionidae – sheathbills (2 species)
Family Pluvianidae – Egyptian plover
Family Charadriidae – plovers (69 species)
Family Recurvirostridae – stilts, avocets (10 species)
Family Ibidorhynchidae – ibisbill
Family Haematopodidae – oystercatchers (12 species)
Suborder Scolopaci
Family Rostratulidae – painted-snipes (3 species)
Family Jacanidae – jacanas (8 species)
Family Pedionomidae – plains-wanderer
Family Thinocoridae – seedsnipes (4 species)
Family Scolopacidae – sandpipers, snipes (98 species)
Suborder Lari
Family Turnicidae – buttonquails (18 species)
Family Dromadidae – crab-plover
Family Glareolidae – coursers, pratincoles (17 species)
Family Laridae – gulls, terns, skimmers (103 species)
Family Stercorariidae – skuas (7 species)
Family Alcidae – auks (25 species)
Evolutionary history
That the Charadriiformes are an ancient group is also borne out by the fossil record. Alongside the Anseriformes, the Charadriiformes are the only other order of modern bird to have an established fossil record within the late Cretaceous, alongside the other dinosaurs. Much of the Neornithes' fossil record around the Cretaceous–Paleogene extinction event is made up of bits and pieces of birds which resemble this order. In many, this is probably due to convergent evolution brought about by semiaquatic habits. Specimen VI 9901 (López de Bertodano Formation, Late Cretaceous of Vega Island, Antarctica) is probably a basal charadriiform somewhat reminiscent of a thick-knee. However, more complete remains of undisputed charadriiforms are known only from the mid-Paleogene onwards. Present-day orders emerged around the Eocene-Oligocene boundary, roughly 35–30 mya. Basal or unresolved charadriiforms are:
"Morsoravis" (Late Paleocene/Early Eocene of Jutland, Denmark) - a nomen nudum?
Jiliniornis (Huadian Middle Eocene of Huadian, China) - charadriid?
Boutersemia (Early Oligocene of Boutersem, Belgium) - glareolid?
Turnipax (Early Oligocene) - turnicid?
Elorius (Early Miocene Saint-Gérand-le-Puy, France)
"Larus" desnoyersii (Early Miocene of SE France) - larid? stercorarid?
"Larus" pristinus (John Day Early Miocene of Willow Creek, US) - larid?
Charadriiformes gen. et sp. indet. (Bathans Early/Middle Miocene of Otago, New Zealand) - charadriid? scolopacid?
Charadriiformes gen. et sp. indet. (Bathans Early/Middle Miocene of Otago, New Zealand) - charadriid? scolopacid?
Charadriiformes gen. et sp. indet. (Bathans Early/Middle Miocene of Otago, New Zealand) - larid?
Charadriiformes gen. et sp. indet. (Sajóvölgyi Middle Miocene of Mátraszõlõs, Hungary
"Totanus" teruelensis (Late Miocene of Los Mansuetos, Spain) - scolopacid? larid?
The "transitional shorebirds" ("Graculavidae") are a generally Mesozoic form taxon formerly believed to constitute the common ancestors of charadriiforms, waterfowl and flamingos. They are now assumed to be mostly basal taxa of the charadriiforms and/or "higher waterbirds", which probably were two distinct lineages 65 mya already, and few if any are still believed to be related to the well-distinct waterfowl. Taxa formerly considered graculavids are:
Laornithidae - charadriiform? gruiform?
Laornis (Late Cretaceous?)
"Graculavidae"
Graculavus (Lance Creek Late Cretaceous - Hornerstown Late Cretaceous/Early Palaeocene) - charadriiform?
Palaeotringa (Hornerstown Late Cretaceous?) - charadriiform?
Telmatornis (Navesink Late Cretaceous?) - charadriiform? gruiform?
Scaniornis - phoenicopteriform?
Zhylgaia - presbyornithid?
Dakotornis
"Graculavidae" gen. et sp. indet. (Gloucester County, US)
Other wader- or gull-like birds incertae sedis, which may or may not be Charadriiformes, are:
Ceramornis (Lance Creek Late Cretaceous)
"Cimolopteryx" (Lance Creek Late Cretaceous)
Palintropus (Lance Creek Late Cretaceous)
Torotix (Late Cretaceous)
Volgavis (Early Paleocene of Volgograd, Russia)
Eupterornis (Paleocene of France)
Neornithes incerta sedis (Late Paleocene/Early Eocene of Ouled Abdoun Basin, Morocco)
Fluviatitavis (Early Eocene of Silveirinha, Portugal)
Evolution of parental care in Charadriiformes
Shorebirds pursue a larger diversity of parental care strategies than do most other avian orders. They therefore present an attractive set of examples to support the understanding of the evolution of parental care in avians generally. The ancestral avian most likely had a female parental care system. The shorebird ancestor specifically evolved from a bi-parental care system, yet the species within the clade Scolopacidae evolved from a male parental care system. These transitions might have occurred for several reasons. Brooding density is correlated with male parental care. Male care systems in birds are shown to have a very low breeding density while female care systems in birds have a high breeding density. (Owens 2005). Certain rates of male and female mortality, male and female egg maturation rate, and egg death rate have been associated with particular systems as well. It has also been shown that sex role reversal is motivated by the male-biased adult sex ratio. The reason for such diversity in shorebirds, compared to other birds, has yet to be understood.
| Biology and health sciences | Charadriiformes | Animals |
53560 | https://en.wikipedia.org/wiki/Procellariiformes | Procellariiformes | Procellariiformes is an order of seabirds that comprises four families: the albatrosses, the petrels and shearwaters, and two families of storm petrels. Formerly called Tubinares and still called tubenoses in English, procellariiforms are often referred to collectively as the petrels, a term that has been applied to all members of the order, or more commonly all the families except the albatrosses. They are almost exclusively pelagic (feeding in the open ocean), and have a cosmopolitan distribution across the world's oceans, with the highest diversity being around New Zealand.
Procellariiforms are colonial, mostly nesting on remote, predator-free islands. The larger species nest on the surface, while most smaller species nest in natural cavities and burrows. They exhibit strong philopatry, returning to their natal colony to breed and returning to the same nesting site over many years. Procellariiforms are monogamous and form long-term pair bonds that are formed over several years and may last for the life of the pair. A single egg is laid per nesting attempt, and usually a single nesting attempt is made per year, although the larger albatrosses may only nest once every two years. Both parents participate in incubation and chick rearing. Incubation times are long compared to other birds, as are fledging periods. Once a chick has fledged there is no further parental care.
Procellariiforms have had a long relationship with humans. They have been important food sources for many people, and continue to be hunted as such in some parts of the world. The albatrosses in particular have been the subject of numerous cultural depictions. Procellariiforms include some of the most endangered bird taxa, with many species threatened with extinction due to introduced predators in their breeding colonies, marine pollution and the danger of fisheries by-catch. Scientists, conservationists, fishermen, and governments around the world are working to reduce the threats posed to them, and these efforts have led to the signing of the Agreement on the Conservation of Albatrosses and Petrels, a legally binding international treaty signed in 2001.
Taxonomy
The order was named Procellariiformes by German anatomist Max Fürbringer in 1888. The word comes from the Latin word procella, which means a violent wind or a storm, and -iformes for order.
Until the beginning of the 20th century, the family Hydrobatidae was named Procellariidae, and the family now called Procellariidae was rendered "Puffinidae." The order itself was called Tubinares. A major early work on this group is Frederick DuCane Godman's Monograph of the Petrels, five fascicles, 1907–1910, with figures by John Gerrard Keulemans.
In the Sibley-Ahlquist taxonomy, the tubenoses were included in a greatly enlarged order "Ciconiiformes". This taxonomic treatment was almost certainly erroneous, but its assumption of a close evolutionary relationship with other "higher waterbirds" – such as loons (Gaviiformes) and penguins (Sphenisciformes) – appears to be correct. The procellariiforms are most closely related to penguins, having diverged from them about 60 million years ago.
The diving petrels in the genus Pelecanoides were formerly placed in their own family Pelecanoididae. When genetic studies found that they were embedded within the family Procellariidae, the two families were merged.
All the storm petrels were once placed in the family Hydrobatidae but genetic data indicated that Hydrobatidae consisted of two deeply divergent clades that were not sister taxa. In 2018 the austral storm petrels were moved to the new family Oceanitidae. The northern storm petrels in the family Hydrobatidae are more closely related to the family Procellariidae than they are to the austral storm petrels in the family Oceanitidae.
Earlier molecular phylogenetic studies found the family Oceantidae containing the austral storm petrels as the most basal with differing branching topologies for other three families. More recent large-scale studies have found a consistent pattern with the albatross family Diomedeidae as the most basal and Hydrobatidae sister to Procellariidae.
There are 147 living species of procellariiform worldwide, and the order is divided into four extant families, with a fifth prehistorically extinct:
Family Diomedeidae (albatrosses) are very large seabirds with a large strong hooked bill. They have strong legs, enabling them to walk well on land.
Family Oceanitidae (Austral storm petrels) are among the smallest seabirds, with fluttering flight and long but weak legs. Most have dark upperparts and a white underside.
Family Hydrobatidae (northern storm petrels) are similar to the austral storm petrels but have longer more pointed wings and most species have forked tails.
Family Procellariidae (shearwaters, fulmarine petrels, gadfly petrels, and prions) are a varied group of small or medium-sized seabirds, the largest being the giant petrels. They are heavy for their size, with a high wing loading, so they need to fly fast. Most, except the giant petrels, have weak legs and are nearly helpless on land.
Family †Diomedeoididae (Early Oligocene – Early Miocene) is an extinct group that had narrow beaks and feet with wide, flat phalanges, especially on the fourth toe.
Fossils of a bird similar to a petrel from the Eocene have been found in the London Clay and in Louisiana. Diving petrels occurred in the Miocene, with a species from that family (Pelecanoides miokuaka) being described in 2007. The most numerous fossils from the Paleogene are those from the extinct family Diomedeoididae, fossils of which have been found in Central Europe and Iran.
Biology
Distribution and movements
The procellariiforms have a cosmopolitan distribution across the world's oceans and seas, although at the levels of family and genus there are some clear patterns. Antarctic petrels, Thalassoica antarctica, have to fly over to get to the ocean from their breeding colonies in Antarctica, and northern fulmars breed on the northeastern tip of Greenland, the northernmost piece of land. The most cosmopolitan family is the Procellariidae, which are found in tropical, temperate and polar zones of both the Northern and the Southern Hemispheres, though the majority do not breed in the tropics, and half the species are restricted to southern temperate and polar regions. The gadfly petrels, Pterodroma, have a generally tropical and temperate distribution, whereas the fulmarine petrels are mostly polar with some temperate species. The majority of the fulmarine petrels, along with the prions, are confined to the Southern Hemisphere.
The storm petrels are almost as widespread as the procellariids, and fall into two distinct families; the Oceanitidae have a mostly Southern Hemisphere distribution and the Hydrobatidae are found mostly in the Northern Hemisphere. Amongst the albatrosses the majority of the family is restricted to the Southern Hemisphere, feeding and nesting in cool temperate areas, although one genus, Phoebastria, ranges across the north Pacific. The family is absent from the north Atlantic, although fossil records indicate they bred there once. Finally the diving petrels are restricted to the Southern Hemisphere.
Migration
The various species within the order have a variety of migration strategies. Some species undertake regular trans-equatorial migrations, such as the sooty shearwater which annually migrates from its breeding grounds in New Zealand and Chile to the North Pacific off Japan, Alaska and California, an annual round trip of , the second longest measured annual migration of any bird. A number of other petrel species undertake trans-equatorial migrations, including the Wilson's storm petrel and the Providence petrel, but no albatrosses cross the equator, as they rely on wind assisted flight. There are other long-distance migrants within the order; Swinhoe's storm petrels breed in the western Pacific and migrate to the western Indian Ocean, and Bonin petrels nesting in Hawaii migrate to the coast of Japan during the non-breeding season.
Navigation
Many species in the order travel long distances over open water but return to the same nest site each year, raising the question of how they navigate so accurately. The Welsh naturalist Ronald Lockley carried out early research into animal navigation with the Manx shearwaters that nested on the island of Skokholm. In release experiments, a Manx shearwater flew from Boston to Skokholm, a distance of in 12 days.
Lockley showed that when released "under a clear sky" with sun or stars visible, the shearwaters oriented themselves and then "flew off in a direct line for Skokholm", making the journey so rapidly that they must have flown almost in a straight line. But if the sky was overcast at the time of release, the shearwaters flew around in circles "as if lost" and returned slowly or not at all, implying that they navigated using astronomical cues.
Researchers have also begun investigating olfaction's role in procellariiform navigation. In a study where Cory's shearwaters were rendered anosmic with zinc sulphate, a compound which kills the surface layer of the olfactory epithelium, and released hundreds of kilometers away from their home colony at night, control birds found their way to their home nests before night was over, whereas anosmic birds did not home until the next day. A similar study that released Cory's shearwaters 800 km from their home nests, testing both magnetic and olfactory disturbances’ effects on navigation, found that anosmic birds took longer to home than magnetically disturbed or control birds.
Morphology and flight
Procellariiforms range in size from the very large wandering albatross, at and a wingspan, to tiny birds like the least storm petrel, at with a wingspan, and the smallest of the prions, the fairy prion, with a wingspan of . Their nostrils are enclosed in one or two tubes on their straight deeply-grooved bills with hooked tips. The beaks are made up of several plates. Their wings are long and narrow; the feet are webbed, and the hind toe is undeveloped or non-existent; their adult plumage is predominantly black, white, and grey.
The order has a few unifying characteristics, starting with their tubular nasal passage which is used for olfaction. Procellariiformes that nest in burrows have a strong sense of smell, being able to detect dimethyl sulfide released from plankton in the ocean. This ability to smell helps to locate patchily distributed prey at sea and may also help locate their nests within nesting colonies. In contrast, surface nesting Procellariiformes have increased vision, having six times better spatial resolution than those that nest in burrows. The structure of the bill, which contains seven to nine distinct horny plates, is another unifying feature, although there are differences within the order. Petrels have a plate called the maxillary unguis that forms a hook on the maxilla. The smaller members of the order have a comb-like mandible, made by the tomial plate, for plankton feeding. Most members of the order are unable to walk well on land, and many species visit their remote breeding islands only at night. The exceptions are the huge albatrosses, several of the gadfly petrels and shearwaters and the fulmar-petrels. The latter can disable even large predatory birds with their obnoxious stomach oil, which they can project some distance. This stomach oil, stored in the proventriculus, is a digestive residue created in the foregut of all tubenoses except the diving petrels, and is used mainly for storage of energy-rich food during their long flights. The oil is also fed to their young, as well as being used for defense.
Procellariiforms drink seawater, so they have to excrete excess salt. All birds have an enlarged nasal gland at the base of the bill, above the eyes, and in the Procellariiformes the gland is active. In general terms, the salt gland removes salt from the system and forms a 5 percent saline solution that drips out of the nostrils, or is forcibly ejected in some petrels. The processes behind this involve high levels of sodium ion reabsorption into the blood plasma within the kidneys, and secretion of sodium chloride via the salt glands using less water than was absorbed, which essentially generates salt-free water for other physiological uses. This high efficiency of sodium ion absorption is attributed to mammalian-type nephrons.
Most albatrosses and procellariids use two techniques to minimise exertion while flying, namely, dynamic soaring and slope soaring. The albatrosses and giant petrels share a morphological adaptation to aid in flight, a sheet of tendon which locks the wing when fully extended, allowing the wing to be kept up and out without any muscle effort. Amongst the Oceanitinae storm-petrels there are two unique flight patterns, one being surface pattering. In this they move across the water surface holding and moving their feet on the water's surface while holding steady above the water, and remaining stationary by hovering with rapid fluttering or by using the wind to anchor themselves in place. A similar flight method is thought to have been used by the extinct petrel family Diomedeoididae. The white-faced storm petrel possesses a unique variation on pattering: holding its wings motionless and at an angle into the wind, it pushes itself off the water's surface in a succession of bounding jumps.
Diet and feeding
The procellariiforms are for the most part exclusively marine foragers; the only exception to this rule are the two species of giant petrel, which regularly feed on carrion or other seabirds while on land. While some other species of fulmarine and Procellaria petrels also take carrion, the diet of most species of albatrosses and petrels is dominated by fish, squid, krill and other marine zooplankton. The importance of these food sources varies from species to species and family to family. For example, of the two albatross species found in Hawaii, the black-footed albatross takes mostly fish, while the Laysan feeds mainly on squid. The albatrosses in general feed on fish, squid and krill. Among the procellariids, the prions concentrate on small crustacea, the fulmarine petrels take fish and krill but little squid, while the Procellaria petrels consume mainly squid. The storm petrels take small droplets of oil from the surface of the water, as well as small crustaceans and fish.
Petrels obtain food by snatching prey while swimming on the surface, snatching prey from the wing or diving down under the water to pursue prey. Dipping down from flight is most commonly used by the gadfly petrels and the storm petrels. There have been records of wedge-tailed shearwaters snatching flying fish from the air, but as a rule this technique is rare. Some diving birds may aid diving by beginning with a plunge from the air, but for the most part petrels are active divers and use their wings to move around under the water. The depths achieved by various species were determined in the 1990s and came as a surprise to scientists; short-tailed shearwaters have been recorded diving to and the Light-mantled sooty albatross to .
Breeding behaviour
Breeding colonies
All procellariiforms are colonial, predominantly breeding on offshore or oceanic islands. The few species that nest on continents do so in inhospitable environments such as dry deserts or on Antarctica. These colonies can vary from the widely spaced colonies of the giant petrels to the dense 3.6 million-strong colonies of Leach's storm petrels. For almost all species the need to breed is the only reason that procellariiforms return to land at all. Some of the larger petrels have to nest on windswept locations as they require wind to take off and forage for food. Within the colonies, pairs defend usually small territories (the giant petrels and some albatrosses can have very large territories) which is the small area around either the nest or a burrow. Competition between pairs can be intense, as is competition between species, particularly for burrows. Larger species of petrels will even kill the chicks and even adults of smaller species in disputes over burrows. Burrows and natural crevices are most commonly used by the smaller species; all the storm petrels and diving petrels are cavity nesters, as are many of the procellariids. The fulmarine petrels and some tropical gadfly petrels and shearwaters are surface nesters, as are all the albatrosses.
Procellariiforms show high levels of philopatry, both site fidelity and natal philopatry. Natal philopatry is the tendency of an individual bird to return to its natal colony to breed, often many years after leaving the colony as a chick. This tendency has been shown through ringing studies and mitochondrial DNA studies. Birds ringed as chicks have been recaptured close to their original nests, sometimes extremely close; in the Laysan albatross the average distance between hatching site and the site where a bird established its own territory was , and a study of Cory's shearwaters nesting near Corsica found that nine out of 61 male chicks that returned to breed at their natal colony actually bred in the burrow they were raised in. Mitochondrial DNA provides evidence of restricted gene flow between different colonies, strongly suggesting philopatry.
The other type of philopatry exhibited is site fidelity, where pairs of birds return to the same nesting site for a number of years. Among the most extreme examples known of this tendency was the fidelity of a ringed northern fulmar that returned to the same nest site for 25 years. The average number of birds returning to the same nest sites is high in all species studied, with around 91 percent for Bulwer's petrels, and 85 percent of males and 76 percent of females for Cory's shearwaters (after a successful breeding attempt).
Pair bonds and life history
Procellariiforms are monogamous breeders and form long-term pair bonds. These pair bonds take several years to develop in some species, particularly with the albatrosses. Once formed, they last for many breeding seasons, in some cases for the life of the pair. Petrel courtship can be elaborate. It reaches its extreme with the albatrosses, where pairs spend many years perfecting and elaborating mating dances. These dances are composed of synchronised performances of various actions such as preening, pointing, calling, bill clacking, staring, and combinations of such behaviours (like the sky-call). Each particular pair will develop their own individual version of the dance. The breeding behaviour of other procellariiforms is less elaborate, although similar bonding behaviours are involved, particularly for surface-nesting species. These can involve synchronised flights, mutual preening and calling. Calls are important for helping birds locate potential mates and distinguishing between species, and may also help individuals assess the quality of potential mates. After pairs have been formed, calls serve to help them reunite; the ability of individuals to recognise their own mate has been demonstrated in several species.
Procellariiforms are K-selected, being long-lived and caring extensively for their few offspring. Breeding is delayed for several years after fledging, sometimes for as long as ten years in the largest species. Once they begin breeding, they make only a single breeding attempt per nesting season; even if the egg is lost early in the season, they seldom re-lay. Much effort is placed into laying a single (proportionally) large egg and raising a single chick. Procellariiforms are long-lived: the longest living albatross known survived for 51 years, but was probably older, and even the tiny storm-petrels are known to have survived for 30 years. Additionally, the oldest living bird is Wisdom, a female Laysan albatross.
Nesting and chick rearing
The majority of procellariiforms nest once a year and do so seasonally. Some tropical shearwaters, like the Christmas shearwater, are able to nest on cycles slightly shorter than a year, and the large great albatrosses (genus Diomedea) nest in alternate years (if successful). Most temperate and polar species nest over the spring-summer, although some albatrosses and procellariids nest over the winter. In the tropics, some species can be found breeding throughout the year, but most nest in discreet periods. Procellariiforms return to nesting colonies as much as several months before laying, and attend their nest sites regularly before copulation. Prior to laying, females embark on a lengthy pre-laying exodus to build up energy reserves in order to lay the exceptionally large egg. In the stormy petrel, a very small procellariiform, the egg can be 29 percent of the body weight of the female, while in the grey-faced petrel, the female may spend as much as 80 days feeding out at sea after courtship before laying the egg.
When the female returns and lays, incubation is shared between the sexes, with the male taking the first incubation stint and the female returning to sea. The duration of individual stints varies from just a few days to as much as several weeks, during which the incubating bird can lose a considerable amount of weight. The incubation period varies from species to species, around 40 days for the smallest storm-petrels but longer for the largest species; for albatrosses it can span 70 to 80 days, which is the longest incubation period of any bird.
Upon hatching, the chicks are semi-precocial, having open eyes, a dense covering of white or grey down feathers, and the ability to move around the nesting site. After hatching, the incubating adult remains with the chick for a number of days, a period known as the guard phase. In the case of most burrow-nesting species, this is only until the chick is able to thermoregulate, usually two or three days. Diving-petrel chicks take longer to thermoregulate and have a longer guard phase than other burrow nesters. However, surface-nesting species, which have to deal with a greater range of weather and to contend with predators like skuas and frigatebirds, consequently have a longer guard phase (as long as two weeks in procellariids and three weeks in albatrosses).
The chick is fed by both parents. Chicks are fed on fish, squid, krill, and stomach oil. Stomach oil is oil composed of neutral dietary lipids that are the residue created by digestion of the prey items. As an energy source for chicks it has several advantages over undigested prey, its calorific value is around 9.6 kcal per gram, which is only slightly lower than the value for diesel oil. This can be a real advantage for species that range over huge distances to provide food for hungry chicks. The oil is also used in defence. All procellariiforms create stomach oil except the diving-petrels.
The chick fledges between two and nine months after hatching, almost twice as long as a gull of the same body mass. The reasons behind the length of time are associated with the distance from the breeding site to food. First, there are few predators at the nesting colonies, therefore there is no pressure to fledge quickly. Second, the time between feedings is long due to the distance from the nest site that adults forage, thus a chick that had a higher growth rate would stand a better chance of starving to death. The duration between feedings vary among species and during the stages of development. Small feeds are frequent during the guard phase, but afterward become less frequent. However, each feed can deliver a large amount of energy; both sooty shearwater and mottled petrel chicks have been recorded to double their weight in a single night, probably when fed by both parents.
Relationship with humans
Role in culture
The most important family culturally is the albatrosses, which have been described by one author as "the most legendary of birds". Albatrosses have featured in poetry in the form of Samuel Taylor Coleridge's famous 1798 poem The Rime of the Ancient Mariner, which in turn gave rise to the usage of albatross as metaphor for a burden. More generally, albatrosses were believed to be good omens, and to kill one would bring bad luck. There are few instances of petrels in culture, although there are sailors' legends regarding the storm petrels, which are considered to warn of oncoming storms. In general, petrels were considered to be "soul birds", representing the souls of drowned sailors, and it was considered unlucky to touch them.
In the Russian language, many petrel species from the Hydrobatidae and Procellariidae families of the order Procellariiformes are known as burevestnik, which literally means 'the announcer of the storm'. When in 1901, the Russian writer Maxim Gorky turned to the imagery of subantarctic avifauna to describe Russian society's attitudes to the coming revolution, he used a storm-announcing petrel as the lead character of a poem that soon became popular in the revolutionary circles as "the battle anthem of the revolution". Although the species called "stormy petrel" in English is not one of those to which the burevestnik name is applied in Russian (it, in fact, is known in Russian as an entirely un-romantic kachurka), the English translators uniformly used the "stormy petrel" image in their translations of the poem, usually known in English as The Song of the Stormy Petrel.
Various tubenose birds are relevant to the mythologies and oral traditions of Polynesia. The Māori used the wing bones of the albatross to carve flutes. In Hawaiian mythology, Laysan albatrosses are considered aumakua, being a sacred manifestation of the ancestors, and quite possibly also the sacred bird of Kāne. The storm petrel features prominently in the "Origin of Birds" myth.
Exploitation
Albatrosses and petrels have been important food sources for humans for as long as people have been able to reach their remote breeding colonies. Amongst the earliest-known examples of this is the remains of shearwaters and albatrosses along with those of other seabirds in 5,000-year-old middens in Chile, although it is likely that they were exploited prior to this. Since then, many other marine cultures, both subsistence and industrial, have exploited procellariiforms, in some cases almost to extinction. Some cultures continue to harvest shearwaters (a practice known as muttonbirding); for example, the Māori of New Zealand use a sustainable traditional method known as kaitiakitanga. In Alaska, residents of Kodiak Island harpoon short-tailed albatrosses, Diomedea albatrus, and until the late 1980s residents of Tristan Island in the Indian Ocean harvested the eggs of the Yellow-nosed Mollymawks, Diomedea chlororhynchos, and sooty albatrosses, Phoebetria fusca. Albatrosses and petrels are also now tourist draws in some locations, such as Taiaroa Head. While such exploitation is non-consumptive, it can have deleterious effects that need careful management to protect both the birds and the tourism.
The English naturalist William Yarrell wrote in 1843 that "ten or twelve years ago, Mr. Gould exhibited twenty-four [storm petrels], in a large dish, at one of the evening meetings of the Zoological Society".
The engraver Thomas Bewick wrote in 1804 that "Pennant, speaking of those [birds] which breed on, or inhabit, the Isle of St Kilda, says—'No bird is of so much use to the islanders as this: the Fulmar supplies them with oil for their lamps, down for their beds, a delicacy for their tables, a balm for their wounds, and a medicine for their distempers.'" A photograph by George Washington Wilson taken about 1886 shows a "view of the men and women of St Kilda on the beach dividing up the catch of Fulmar". James Fisher, author of The Fulmar (1952) calculated that every person on St Kilda consumed over 100 fulmars each year; the meat was their staple food, and they caught around 12,000 birds annually. However, when the human population left St Kilda in 1930, the population did not suddenly grow.
Threats and conservation
The albatrosses and petrels are "amongst the most severely threatened taxa worldwide". They face a variety of threats, the severity of which varies greatly from species to species. Several species are among the most common of seabirds, including Wilson's storm petrel (an estimated 12 to 30 million individuals) and the short-tailed shearwater (23 million individuals); while the total population of some other species is a few hundred. There are less than 200 Magenta petrels breeding on the Chatham Islands, only 130 to 160 Zino's petrels and only 170 Amsterdam albatrosses. Only one species is thought to have become extinct since 1600, the Guadalupe storm petrel of Mexico, although a number of species had died out before this. Numerous species are very poorly known; for example, the Fiji petrel has rarely been seen since its discovery. The breeding colony of the New Zealand storm petrel was not located until February 2013; it had been thought extinct for 150 years until its rediscovery in 2003, while the Bermuda petrel had been considered extinct for nearly 300 years.
The principal threat to the albatrosses and larger species of procellariids is long-line fishing. Bait set on hooks is attractive to foraging birds and many are hooked by the lines as they are set. As many as 100,000 albatrosses are hooked and drown each year on tuna lines set out by long-line fisheries. Before 1991 and the ban on drift-net fisheries, it was estimated that 500,000 seabirds a year died as a result. This has caused steep declines in some species, as procellariiforms are extremely slow breeders and cannot replace their numbers fast enough. Losses of albatrosses and petrels in the Southern Ocean were estimated at between 1 percent and 16 percent per year, which these species cannot sustain for long.
Exotic species introduced to the remote breeding colonies threaten all types of procellariiform. These principally take the form of predators; most albatross and petrel species are clumsy on land and unable to defend themselves from mammals such as rats, feral cats and pigs. This phenomenon, ecological naivete, has resulted in declines in many species and was implicated in the extinction of the Guadalupe storm petrel. Already in 1910 Godman wrote:
Introduced herbivores may unbalance the ecology of islands; introduced rabbits destroyed the forest understory on Cabbage Tree Island off New South Wales, which increased the vulnerability of the Gould's petrels nesting on the island to natural predators, and left them vulnerable to the sticky fruits of the native birdlime tree (Pisonia umbellifera). In the natural state these fruits lodge in the understory of the forest, but with the understory removed the fruits fall to the ground where the petrels move about, sticking to their feathers and making flight impossible.
Exploitation has decreased in importance as a threat. Other threats include the ingestion of plastic flotsam. Once swallowed, plastic can cause a general decline in the fitness of the bird, or in some cases lodge in the gut and cause a blockage, leading to death by starvation. It can also be picked up by foraging adults and fed to chicks, stunting their development and reducing the chances of successfully fledging. Procellariids are also vulnerable to marine pollution, as well as oil spills. Some species, such as Barau's petrel, Newell's shearwater and Cory's shearwater, which nest high up on large developed islands, are victims of light pollution. Fledging chicks are attracted to streetlights and may then be unable to reach the sea. An estimated 20 to 40 percent of fledging Barau's petrels and 45 to 60 percent of fledging Cory's shearwater are attracted to the streetlights on Réunion and Tenerife, respectively.
| Biology and health sciences | Procellariiformes | null |
53601 | https://en.wikipedia.org/wiki/Metal%20casting | Metal casting | In metalworking and jewelry making, casting is a process in which a liquid metal is delivered into a mold (usually by a crucible) that contains a negative impression (i.e., a three-dimensional negative image) of the intended shape. The metal is poured into the mold through a hollow channel called a sprue. The metal and mold are then cooled, and the metal part (the casting) is extracted. Casting is most often used for making complex shapes that would be difficult or uneconomical to make by other methods.
Casting processes have been known for thousands of years, and have been widely used for sculpture (especially in bronze), jewelry in precious metals, and weapons and tools. Highly engineered castings are found in 90 percent of durable goods, including cars, trucks, aerospace, trains, mining and construction equipment, oil wells, appliances, pipes, hydrants, wind turbines, nuclear plants, medical devices, defense products, toys, and more.
Traditional techniques include lost-wax casting (which may be further divided into centrifugal casting, and vacuum assist direct pour casting), plaster mold casting and sand casting.
The modern casting process is subdivided into two main categories: expendable and non-expendable casting. It is further broken down by the mold material, such as sand or metal, and pouring method, such as gravity, vacuum, or low pressure.
Expendable mold casting
Expendable mold casting is a generic classification that includes sand, plastic, shell, plaster, and investment (lost-wax technique) moldings. This method of mold casting involves the use of temporary, non-reusable molds.
Sand casting
Sand casting is one of the most popular and simplest types of casting, and has been used for centuries. Sand casting allows for smaller batches than permanent mold casting and at a very reasonable cost. Not only does this method allow manufacturers to create products at a low cost, but there are other benefits to sand casting, such as very small-size operations. The process allows for castings small enough fit in the palm of one's hand to those large enough for a train car bed (one casting can create the entire bed for one rail car). Sand casting also allows most metals to be cast depending on the type of sand used for the molds.
Sand casting requires a lead time of days, or even weeks sometimes, for production at high output rates (1–20 pieces/hr-mold) and is unsurpassed for large-part production. Green (moist) sand, which is black in color, has almost no part weight limit, whereas dry sand has a practical part mass limit of . Minimum part weight ranges from . The sand is bonded using clays, chemical binders, or polymerized oils (such as motor oil). Sand can be recycled many times in most operations and requires little maintenance.
Loam molding
Loam molding has been used to produce large symmetrical objects such as cannon and church bells. Loam is a mixture of clay and sand with straw or dung. A model of the produced is formed in a friable material (the chemise). The mold is formed around this chemise by covering it with loam. This is then baked (fired) and the chemise removed. The mold is then stood upright in a pit in front of the furnace for the molten metal to be poured. Afterwards the mold is broken off. Molds can thus only be used once, so that other methods are preferred for most purposes.
Plaster mold casting
Plaster casting is similar to sand casting except that plaster of paris is used instead of sand as a mold material. Generally, the form takes less than a week to prepare, after which a production rate of 1–10 units/hr-mold is achieved, with items as massive as and as small as with very good surface finish and close tolerances. Plaster casting is an inexpensive alternative to other molding processes for complex parts due to the low cost of the plaster and its ability to produce near net shape castings. The biggest disadvantage is that it can only be used with low melting point non-ferrous materials, such as aluminium, copper, magnesium, and zinc.
Shell molding
Shell molding is similar to sand casting, but the molding cavity is formed by a hardened "shell" of sand instead of a flask filled with sand. The sand used is finer than sand casting sand and is mixed with a resin so that it can be heated by the pattern and hardened into a shell around the pattern. Because of the resin and finer sand, it gives a much finer surface finish. The process is easily automated and more precise than sand casting. Common metals that are cast include cast iron, aluminium, magnesium, and copper alloys. This process is ideal for complex items that are small to medium-sized.
Investment casting
Investment casting (known as lost-wax casting in art) is a process that has been practiced for thousands of years, with the lost-wax process being one of the oldest known metal forming techniques. From 5000 years ago, when beeswax formed the pattern, to today's high technology waxes, refractory materials, and specialist alloys, the castings ensure high-quality components are produced with the key benefits of accuracy, repeatability, versatility, and integrity.
Investment casting derives its name from the fact that the pattern is invested, or surrounded, with a refractory material. The wax patterns require extreme care for they are not strong enough to withstand forces encountered during the mold making. One advantage of investment casting is that the wax can be reused.
The process is suitable for repeatable production of net shape components from a variety of different metals and high performance alloys. Although generally used for small castings, this process has been used to produce complete aircraft door frames, with steel castings of up to 300 kg and aluminium castings of up to 30 kg. Compared to other casting processes such as die casting or sand casting, it can be an expensive process. However, the components that can be produced using investment casting can incorporate intricate contours, and in most cases the components are cast near net shape, so require little or no rework once cast.
Waste molding of plaster
A durable plaster intermediate is often used as a stage toward the production of a bronze sculpture or as a pointing guide for the creation of a carved stone. With the completion of a plaster, the work is more durable (if stored indoors) than a clay original which must be kept moist to avoid cracking. With the low cost plaster at hand, the expensive work of bronze casting or stone carving may be deferred until a patron is found, and as such work is considered to be a technical, rather than artistic process, it may even be deferred beyond the lifetime of the artist.
In waste molding a simple and thin plaster mold, reinforced by sisal or burlap, is cast over the original clay mixture. When cured, it is then removed from the damp clay, incidentally destroying the fine details in undercuts present in the clay, but which are now captured in the mold. The mold may then at any later time (but only once) be used to cast a plaster positive image, identical to the original clay. The surface of this plaster may be further refined and may be painted and waxed to resemble a finished bronze casting.
Evaporative-pattern casting
This is a class of casting processes that use pattern materials that evaporate during the pour, which means there is no need to remove the pattern material from the mold before casting. The two main processes are lost-foam casting and full-mold casting.
Lost-foam casting
Lost-foam casting is a type of evaporative-pattern casting process that is similar to investment casting except foam is used for the pattern instead of wax. This process takes advantage of the low boiling point of foam to simplify the investment casting process by removing the need to melt the wax out of the mold.
Full-mold casting
Full-mold casting is an evaporative-pattern casting process which is a combination of sand casting and lost-foam casting. It uses an expanded polystyrene foam pattern which is then surrounded by sand, much like sand casting. The metal is then poured directly into the mold, which vaporizes the foam upon contact.
Non-expendable mold casting
Non-expendable mold casting differs from expendable processes in that the mold need not be reformed after each production cycle. This technique includes at least four different methods: permanent, die, centrifugal, and continuous casting. This form of casting also results in improved repeatability in parts produced and delivers near net shape results.
Permanent mold casting
Permanent mold casting is a metal casting process that employs reusable molds ("permanent molds"), usually made from metal. The most common process uses gravity to fill the mold. However, gas pressure or a vacuum are also used. A variation on the typical gravity casting process, called slush casting, produces hollow castings. Common casting metals are aluminum, magnesium, and copper alloys. Other materials include tin, zinc, and lead alloys and iron and steel are also cast in graphite molds. Permanent molds, while lasting more than one casting still have a limited life before wearing out.
Die casting
The die casting process forces molten metal under high pressure into mold cavities (which are machined into dies). Most die castings are made from nonferrous metals, specifically zinc, copper, and aluminium-based alloys, but ferrous metal die castings are possible. The die casting method is especially suited for applications where many small to medium-sized parts are needed with good detail, a fine surface quality and dimensional consistency.
Semi-solid metal casting
Semi-solid metal (SSM) casting is a modified die casting process that reduces or eliminates the residual porosity present in most die castings. Rather than using liquid metal as the feed material, SSM casting uses a higher viscosity feed material that is partially solid and partially liquid. A modified die casting machine is used to inject the semi-solid slurry into reusable hardened steel dies. The high viscosity of the semi-solid metal, along with the use of controlled die filling conditions, ensures that the semi-solid metal fills the die in a non-turbulent manner so that harmful porosity can be essentially eliminated.
Used commercially mainly for aluminium and magnesium alloys, SSM castings can be heat treated to the T4, T5 or T6 tempers. The combination of heat treatment, fast cooling rates (from using uncoated steel dies) and minimal porosity provides excellent combinations of strength and ductility. Other advantages of SSM casting include the ability to produce complex shaped parts net shape, pressure tightness, tight dimensional tolerances and the ability to cast thin walls.
Centrifugal casting
In this process molten metal is poured in the mold and allowed to solidify while the mold is rotating. Metal is poured into the center of the mold at its axis of rotation. Due to inertial force, the liquid metal is thrown out toward the periphery.
Centrifugal casting is both gravity and pressure independent since it creates its own force feed using a temporary sand mold held in a spinning chamber. Lead time varies with the application. Semi- and true-centrifugal processing permit 30–50 pieces/hr-mold to be produced, with a practical limit for batch processing of approximately 9000 kg total mass with a typical per-item limit of 2.3–4.5 kg.
Industrially, the centrifugal casting of railway wheels was an early application of the method developed by the German industrial company Krupp and this capability enabled the rapid growth of the enterprise.
Small art pieces such as jewelry are often cast by this method using the lost wax process, as the forces enable the rather viscous liquid metals to flow through very small passages and into fine details such as leaves and petals. This effect is similar to the benefits from vacuum casting, also applied to jewelry casting.
Continuous casting
Continuous casting is a refinement of the casting process for the continuous, high-volume production of metal sections with a constant cross-section. It's primarily used to produce a semi-finished products for further processing. Molten metal is poured into an open-ended, water-cooled mold, which allows a 'skin' of solid metal to form over the still-liquid center, gradually solidifying the metal from the outside in. After solidification, the strand, as it is sometimes called, is continuously withdrawn from the mold. Predetermined lengths of the strand can be cut off by either mechanical shears or traveling oxyacetylene torches and transferred to further forming processes, or to a stockpile. Cast sizes can range from strip (a few millimeters thick by about five meters wide) to billets (90 to 160 mm square) to slabs (1.25 m wide by 230 mm thick). Sometimes, the strand may undergo an initial hot rolling process before being cut.
Continuous casting is used due to the lower costs associated with continuous production of a standard product, and also increased quality of the final product. Metals such as steel, copper, aluminum and lead are continuously cast, with steel being the metal with the greatest tonnages cast using this method.
Upcasting
The upcasting (up-casting, upstream, or upward casting) is a method of either vertical or horizontal continuous casting of rods and pipes of various profiles (cylindrical, square, hexagonal, slabs etc.) of 8-30mm in diameter. Copper (Cu), bronze (Cu·Sn alloy), nickel alloys are usually used because of greater casting speed (in case of vertical upcasting) and because of better physical features obtained. The advantage of this method is that metals are almost oxygen-free and that the rate of product crystallization (solidification) may be adjusted in a crystallizer - a high-temperature resistant device that cools a growing metal rod or pipe by using water.
The method is comparable to Czochralski method of growing silicon (Si) crystals, which is a metalloid.
Terminology
Metal casting processes uses the following terminology:
Pattern: An approximate duplicate of the final casting used to form the mold cavity.
Molding material: The material that is packed around the pattern and then the pattern is removed to leave the cavity where the casting material will be poured.
Flask: The rigid wood or metal frame that holds the molding material.
Cope: The top half of the pattern, flask, mold, or core.
Drag: The bottom half of the pattern, flask, mold, or core.
Core: An insert in the mold that produces internal features in the casting, such as holes.
Core print: The region added to the pattern, core, or mold used to locate and support the core.
Mold cavity: The combined open area of the molding material and core, where the metal is poured to produce the casting.
Riser: An extra void in the mold that fills with molten material to compensate for shrinkage during solidification.
Gating system: The network of connected channels that deliver the molten material to the mold cavities.
Pouring cup or pouring basin: The part of the gating system that receives the molten material from the pouring vessel.
Sprue: The pouring cup attaches to the sprue, which is the vertical part of the gating system. The other end of the sprue attaches to the runners.
Runners: The horizontal portion of the gating system that connects the sprues to the gates.
Gates: The controlled entrances from the runners into the mold cavities.
Vents: Additional channels that provide an escape for gases generated during the pour.
Parting line or parting surface: The interface between the cope and drag halves of the mold, flask, or pattern.
Draft: The taper on the casting or pattern that allow it to be withdrawn from the mold
Core box: The mold or die used to produce the cores.
Chaplet: Long vertical holding rod for core that after casting it become the integral part of casting, provide the support to the core.
Some specialized processes, such as die casting, use additional terminology.
Theory
Casting is a solidification process, which means the solidification phenomenon controls most of the properties of the casting. Moreover, most of the casting defects occur during solidification, such as gas porosity and solidification shrinkage.
Solidification occurs in two steps: nucleation and crystal growth. In the nucleation stage, solid particles form within the liquid. When these particles form, their internal energy is lower than the surrounded liquid, which creates an energy interface between the two. The formation of the surface at this interface requires energy, so as nucleation occurs, the material actually undercools (i.e. cools below its solidification temperature) because of the extra energy required to form the interface surfaces. It then recalescences, or heats back up to its solidification temperature, for the crystal growth stage. Nucleation occurs on a pre-existing solid surface because not as much energy is required for a partial interface surface as for a complete spherical interface surface. This can be advantageous because fine-grained castings possess better properties than coarse-grained castings. A fine grain structure can be induced by grain refinement or inoculation, which is the process of adding impurities to induce nucleation.
All of the nucleations represent a crystal, which grows as the heat of fusion is extracted from the liquid until there is no liquid left. The direction, rate, and type of growth can be controlled to maximize the properties of the casting. Directional solidification is when the material solidifies at one end and proceeds to solidify to the other end; this is the most ideal type of grain growth because it allows liquid material to compensate for shrinkage.
Cooling curves
Cooling curves are important in controlling the quality of a casting. The most important part of the cooling curve is the cooling rate which affects the microstructure and properties. Generally speaking, an area of the casting which is cooled quickly will have a fine grain structure and an area which cools slowly will have a coarse grain structure. Below is an example cooling curve of a pure metal or eutectic alloy, with defining terminology.
Note that before the thermal arrest the material is a liquid and after it the material is a solid; during the thermal arrest the material is converting from a liquid to a solid. Also, note that the greater the superheat the more time there is for the liquid material to flow into intricate details.
The above cooling curve depicts a basic situation with a pure metal, however, most castings are of alloys, which have a cooling curve shaped as shown below.
Note that there is no longer a thermal arrest, instead there is a freezing range. The freezing range corresponds directly to the liquidus and solidus found on the phase diagram for the specific alloy.
Chvorinov's rule
The local solidification time can be calculated using Chvorinov's rule, which is:
Where t is the solidification time, V is the volume of the casting, A is the surface area of the casting that contacts the mold, n is a constant, and B is the mold constant. It is most useful in determining if a riser will solidify before the casting, because if the riser does solidify first then it is worthless.
The gating system
The gating system serves many purposes, the most important being conveying the liquid material to the mold, but also controlling shrinkage, the speed of the liquid, turbulence, and trapping dross. The gates are usually attached to the thickest part of the casting to assist in controlling shrinkage. In especially large castings multiple gates or runners may be required to introduce metal to more than one point in the mold cavity. The speed of the material is important because if the material is traveling too slowly it can cool before completely filling, leading to misruns and cold shuts. If the material is moving too fast then the liquid material can erode the mold and contaminate the final casting. The shape and length of the gating system can also control how quickly the material cools; short round or square channels minimize heat loss.
The gating system may be designed to minimize turbulence, depending on the material being cast. For example, steel, cast iron, and most copper alloys are turbulent insensitive, but aluminium and magnesium alloys are turbulent sensitive. The turbulent insensitive materials usually have a short and open gating system to fill the mold as quickly as possible. However, for turbulent sensitive materials short sprues are used to minimize the distance the material must fall when entering the mold. Rectangular pouring cups and tapered sprues are used to prevent the formation of a vortex as the material flows into the mold; these vortices tend to suck gas and oxides into the mold. A large sprue well is used to dissipate the kinetic energy of the liquid material as it falls down the sprue, decreasing turbulence. The choke, which is the smallest cross-sectional area in the gating system used to control flow, can be placed near the sprue well to slow down and smooth out the flow. Note that on some molds the choke is still placed on the gates to make separation of the part easier, but induces extreme turbulence. The gates are usually attached to the bottom of the casting to minimize turbulence and splashing.
The gating system may also be designed to trap dross. One method is to take advantage of the fact that some dross has a lower density than the base material so it floats to the top of the gating system. Therefore, long flat runners with gates that exit from the bottom of the runners can trap dross in the runners; note that long flat runners will cool the material more rapidly than round or square runners. For materials where the dross is a similar density to the base material, such as aluminium, runner extensions and runner wells can be advantageous. These take advantage of the fact that the dross is usually located at the beginning of the pour, therefore the runner is extended past the last gate(s) and the contaminates are contained in the wells. Screens or filters may also be used to trap contaminates.
It is important to keep the size of the gating system small, because it all must be cut from the casting and remelted to be reused. The efficiency, or , of a casting system can be calculated by dividing the weight of the casting by the weight of the metal poured. Therefore, the higher the number the more efficient the gating system/risers.
Shrinkage
There are three types of shrinkage: shrinkage of the liquid, solidification shrinkage and patternmaker's shrinkage. The shrinkage of the liquid is rarely a problem because more material is flowing into the mold behind it. Solidification shrinkage occurs because metals are less dense as a liquid than a solid, so during solidification the metal density dramatically increases. Patternmaker's shrinkage refers to the shrinkage that occurs when the material is cooled from the solidification temperature to room temperature, which occurs due to thermal contraction.
Solidification shrinkage
Most materials shrink as they solidify, but, as the adjacent table shows, a few materials do not, such as gray cast iron. For the materials that do shrink upon solidification the type of shrinkage depends on how wide the freezing range is for the material. For materials with a narrow freezing range, less than , a cavity, known as a pipe, forms in the center of the casting, because the outer shell freezes first and progressively solidifies to the center. Pure and eutectic metals usually have narrow solidification ranges. These materials tend to form a skin in open air molds, therefore they are known as skin forming alloys. For materials with a wide freezing range, greater than , much more of the casting occupies the mushy or slushy zone (the temperature range between the solidus and the liquidus), which leads to small pockets of liquid trapped throughout and ultimately porosity. These castings tend to have poor ductility, toughness, and fatigue resistance. Moreover, for these types of materials to be fluid-tight, a secondary operation is required to impregnate the casting with a lower melting point metal or resin.
For the materials that have narrow solidification ranges, pipes can be overcome by designing the casting to promote directional solidification, which means the casting freezes first at the point farthest from the gate, then progressively solidifies toward the gate. This allows a continuous feed of liquid material to be present at the point of solidification to compensate for the shrinkage. Note that there is still a shrinkage void where the final material solidifies, but if designed properly, this will be in the gating system or riser.
Risers and riser aids
Risers, also known as feeders, are the most common way of providing directional solidification. It supplies liquid metal to the solidifying casting to compensate for solidification shrinkage. For a riser to work properly the riser must solidify after the casting, otherwise it cannot supply liquid metal to shrinkage within the casting. Risers add cost to the casting because it lowers the yield of each casting; i.e. more metal is lost as scrap for each casting. Another way to promote directional solidification is by adding chills to the mold. A chill is any material which will conduct heat away from the casting more rapidly than the material used for molding.
Risers are classified by three criteria. The first is if the riser is open to the atmosphere, if it is then it is called an open riser, otherwise it is known as a blind type. The second criterion is where the riser is located; if it is located on the casting then it is known as a top riser and if it is located next to the casting it is known as a side riser. Finally, if the riser is located on the gating system so that it fills after the molding cavity, it is known as a live riser or hot riser, but if the riser fills with materials that have already flowed through the molding cavity it is known as a dead riser or cold riser.
Riser aids are items used to assist risers in creating directional solidification or reducing the number of risers required. One of these items are chills which accelerate cooling in a certain part of the mold. There are two types: external and internal chills. External chills are masses of high-heat-capacity and high-thermal-conductivity material that are placed on an edge of the molding cavity. Internal chills are pieces of the same metal that is being poured, which are placed inside the mold cavity and become part of the casting. Insulating sleeves and toppings may also be installed around the riser cavity to slow the solidification of the riser. Heater coils may also be installed around or above the riser cavity to slow solidification.
Patternmaker's shrink
Shrinkage after solidification can be dealt with by using an oversized pattern designed specifically for the alloy used. s, or s, are used to make the patterns oversized to compensate for this type of shrinkage. These rulers are up to 2.5% oversize, depending on the material being cast. These rulers are mainly referred to by their percentage change. A pattern made to match an existing part would be made as follows: First, the existing part would be measured using a standard ruler, then when constructing the pattern, the pattern maker would use a contraction rule, ensuring that the casting would contract to the correct size.
Note that patternmaker's shrinkage does not take phase change transformations into account. For example, eutectic reactions, martensitic reactions, and graphitization can cause expansions or contractions.
Mold cavity
The mold cavity of a casting does not reflect the exact dimensions of the finished part due to a number of reasons. These modifications to the mold cavity are known as allowances and account for patternmaker's shrinkage, draft, machining, and distortion. In non-expendable processes, these allowances are imparted directly into the permanent mold, but in expendable mold processes they are imparted into the patterns, which later form the mold cavity. Note that for non-expendable molds an allowance is required for the dimensional change of the mold due to heating to operating temperatures.
For surfaces of the casting that are perpendicular to the parting line of the mold a draft must be included. This is so that the casting can be released in non-expendable processes or the pattern can be released from the mold without destroying the mold in expendable processes. The required draft angle depends on the size and shape of the feature, the depth of the mold cavity, how the part or pattern is being removed from the mold, the pattern or part material, the mold material, and the process type. Usually the draft is not less than 1%.
The machining allowance varies drastically from one process to another. Sand castings generally have a rough surface finish, therefore need a greater machining allowance, whereas die casting has a very fine surface finish, which may not need any machining tolerance. Also, the draft may provide enough of a machining allowance to begin with.
The distortion allowance is only necessary for certain geometries. For instance, U-shaped castings will tend to distort with the legs splaying outward, because the base of the shape can contract while the legs are constrained by the mold. This can be overcome by designing the mold cavity to slope the leg inward to begin with. Also, long horizontal sections tend to sag in the middle if ribs are not incorporated, so a distortion allowance may be required.
Cores may be used in expendable mold processes to produce internal features. The core can be of metal but it is usually done in sand.
Filling
There are a few common methods for filling the mold cavity: gravity, low-pressure, high-pressure, and vacuum.
Vacuum filling, also known as counter-gravity filling, is more metal efficient than gravity pouring because less material solidifies in the gating system. Gravity pouring only has a 15 to 50% metal yield as compared to 60 to 95% for vacuum pouring. There is also less turbulence, so the gating system can be simplified since it does not have to control turbulence. Plus, because the metal is drawn from below the top of the pool the metal is free from dross and slag, as these are lower density (lighter) and float to the top of the pool. The pressure differential helps the metal flow into every intricacy of the mold. Finally, lower temperatures can be used, which improves the grain structure. The first patented vacuum casting machine and process dates to 1879.
Low-pressure filling uses 5 to 15 psig (35 to 100 kPag) of air pressure to force liquid metal up a feed tube into the mold cavity. This eliminates turbulence found in gravity casting and increases density, repeatability, tolerances, and grain uniformity. After the casting has solidified the pressure is released and any remaining liquid returns to the crucible, which increases yield.
Tilt filling
Tilt filling, also known as tilt casting, is an uncommon filling technique where the crucible is attached to the gating system and both are slowly rotated so that the metal enters the mold cavity with little turbulence. The goal is to reduce porosity and inclusions by limiting turbulence. For most uses tilt filling is not feasible because the following inherent problem: if the system is rotated slow enough to not induce turbulence, the front of the metal stream begins to solidify, which results in mis-runs. If the system is rotated faster it induces turbulence, which defeats the purpose. Durville of France was the first to try tilt casting, in the 1800s. He tried to use it to reduce surface defects when casting coinage from aluminium bronze.
Macrostructure
The grain macrostructure in ingots and most castings have three distinct regions or zones: the chill zone, columnar zone, and equiaxed zone. The image below depicts these zones.
The chill zone is named so because it occurs at the walls of the mold where the wall chills the material. Here is where the nucleation phase of the solidification process takes place. As more heat is removed the grains grow towards the center of the casting. These are thin, long columns that are perpendicular to the casting surface, which are undesirable because they have anisotropic properties. Finally, in the center the equiaxed zone contains spherical, randomly oriented crystals. These are desirable because they have isotropic properties. The creation of this zone can be promoted by using a low pouring temperature, alloy inclusions, or inoculants.
Inspection
Common inspection methods for steel castings are magnetic particle testing and liquid penetrant testing. Common inspection methods for aluminum castings are radiography, ultrasonic testing, and liquid penetrant testing.
Defects
There are a number of problems that can be encountered during the casting process. The main types are: gas porosity, shrinkage defects, mold material defects, pouring metal defects, and metallurgical defects.
Casting process simulation
Casting processes simulation uses numerical methods to calculate cast component quality considering mold filling, solidification and cooling, and provides a quantitative prediction of casting mechanical properties, thermal stresses and distortion. Simulation accurately describes a cast component's quality up-front before production starts. The casting rigging can be designed with respect to the required component properties. This has benefits beyond a reduction in pre-production sampling, as the precise layout of the complete casting system also leads to energy, material, and tooling savings.
The software supports the user in component design, the determination of melting practice and casting methoding through to pattern and mold making, heat treatment, and finishing. This saves costs along the entire casting manufacturing route.
Casting process simulation was initially developed at universities starting from the early 1970s, mainly in Europe and in the U.S., and is regarded as the most important innovation in casting technology over the last 50 years. Since the late 1980s, commercial programs are available which make it possible for foundries to gain new insight into what is happening inside the mold or die during the casting process.
| Technology | Metallurgy | null |
53603 | https://en.wikipedia.org/wiki/Double%20star | Double star | In observational astronomy, a double star or visual double is a pair of stars that appear close to each other as viewed from Earth, especially with the aid of optical telescopes.
This occurs because the pair either forms a binary star (i.e. a binary system of stars in mutual orbit, gravitationally bound to each other) or is an optical double, a chance line-of-sight alignment of two stars at different distances from the observer. Binary stars are important to stellar astronomers as knowledge of their motions allows direct calculation of stellar mass and other stellar parameters. The only (possible) case of "binary star" whose two components are separately visible to the naked eye is the case of Mizar and Alcor (though actually a multiple-star system), but it is not known for certain whether Mizar and Alcor are gravitationally bound.
Since the beginning of the 1780s, both professional and amateur double star observers have telescopically measured the distances and angles between double stars to determine the relative motions of the pairs. If the relative motion of a pair determines a curved arc of an orbit, or if the relative motion is small compared to the common proper motion of both stars, it may be concluded that the pair is in mutual orbit as a binary star. Otherwise, the pair is optical. Multiple stars are also studied in this way, although the dynamics of multiple stellar systems are more complex than those of binary stars.
The following are three types of paired stars:
Optical doubles are unrelated stars that appear close together through chance alignment with Earth.
Visual binaries are gravitationally bound stars that are separately visible with a telescope.
Non-visual binaries are stars whose binary status was deduced through more esoteric means, such as occultation (eclipsing binaries), spectroscopy (spectroscopic binaries), or anomalies in proper motion (astrometric binaries).
Improvements in telescopes can shift previously non-visual binaries into visual binaries, as happened with Polaris A in 2006. It is only the inability to telescopically observe two separate stars that distinguishes non-visual and visual binaries.
History
Mizar, in Ursa Major, was observed to be double by Benedetto Castelli and Galileo. The identification of other doubles soon followed: Robert Hooke discovered one of the first double-star systems, Gamma Arietis, in 1664, while the bright southern star Acrux, in the Southern Cross, was discovered to be double by Fontenay in 1685. Since that time, the search has been carried out thoroughly and the entire sky has been examined for double stars down to a limiting apparent magnitude of about 9.0. At least 1 in 18 stars brighter than 9.0 magnitude in the northern half of the sky are known to be double stars visible with a telescope.
The unrelated categories of optical doubles and true binaries are lumped together for historical and practical reasons. When Mizar was found to be a binary, it was quite difficult to determine whether a double star was a binary system or only an optical double. Improved telescopes, spectroscopy, and photography are the basic tools used to make the distinction. After it was determined to be a visual binary, Mizar's components were found to be spectroscopic binaries themselves.
Observation of double stars
Observation of visual double stars by visual measurement will yield the separation, or angular distance, between the two component stars in the sky and the position angle. The position angle specifies the direction in which the stars are separated and is defined as the bearing from the brighter component to the fainter, where north is 0°. These measurements are called measures. In the measures of a visual binary, the position angle will change progressively and the separation between the two stars will oscillate between maximum and minimum values. Plotting the measures in the plane will produce an ellipse. This is the apparent orbit, the projection of the orbit of the two stars onto the celestial sphere; the true orbit can be computed from it. Although it is expected that the majority of catalogued visual doubles are visual binaries, orbits have been computed for only a few thousand of the over 100,000 known visual double stars.
Distinction between binary stars and other double stars
Confirmation of a visual double star as a binary star can be achieved by observing the relative motion of the components. If the motion is part of an orbit, or if the stars have similar radial velocities or the difference in their proper motions is small compared to their common proper motion, the pair is probably physical. When observed over a short period of time, the components of both optical doubles and long-period visual binaries will appear to be moving in straight lines; for this reason, it can be difficult to distinguish between these two possibilities.
Designations
Some bright visual double stars have a Bayer designation. In this case, the components may be denoted by superscripts. An example of this is α Crucis (Acrux), whose components are α1 Crucis and α2 Crucis. Since α1 Crucis is a spectroscopic binary, this is actually a multiple star. Superscripts are also used to distinguish more distant, physically unrelated, pairs of stars with the same Bayer designation, such as α1,2 Capricorni, ξ1,2 Centauri, and ξ1,2 Sagittarii. These optical pairs are resolvable by the naked eye.
Apart from these pairs, the components of a double star are generally denoted by the letters A (for the brighter, primary, star) and B (for the fainter, secondary, star) appended to the designation, of whatever sort, of the double star. For example, the components of α Canis Majoris (Sirius) are α Canis Majoris A and α Canis Majoris B (Sirius A and Sirius B); the components of 44 Boötis are 44 Boötis A and 44 Boötis B; the components of ADS 16402 are ADS 16402A and ADS 16402B; and so on. The letters AB may be used together to designate the pair. In the case of multiple stars, the letters C, D, and so on may be used to denote additional components, often in order of increasing separation from the brightest star, A.
Visual doubles are also designated by an abbreviation for the name of their discoverer followed by a catalogue number unique to that observer. For example, the pair α Centauri AB was discovered by Father Richaud in 1689, and so is designated RHD 1. Other examples include Δ65, the 65th double discovered by James Dunlop, and Σ2451, discovered by F. G. W. Struve.
The Washington Double Star Catalog, a large database of double and multiple stars, contains over 100,000 entries, each of which gives measures for the separation of two components. Each double star forms one entry in the catalog; multiple stars with n components will be represented by entries in the catalog for n−1 pairs, each giving the separation of one component of the multiple star from another. Codes such as AC are used to denote which components are being measured—in this case, component C relative to component A. This may be altered to a form such as AB-D to indicate the separation of a component from a close pair of components (in this case, component D relative to the pair AB.) Codes such as Aa may also be used to denote a component which is being measured relative to another component, A in this case. Discoverer designations are also listed; however, traditional discoverer abbreviations such as Δ and Σ have been encoded into a string of uppercase Roman letters, so that, for example, Δ65 has become DUN 65 and Σ2451 has become STF 2451. Further examples of this are shown in the adjacent table.
Examples
Visual binaries
Acrux
Capella
p Eridani
Gamma Leonis
Gamma Andromedae
Polaris
Procyon
Sirius
Alpha Centauri system (AB) and Proxima Centauri (thus α Cen C): Actually a three-star system
Optical doubles
Alpha1 and Alpha2 Capricorni
Theta Muscae and Theta Muscae B
Zeta1 and Zeta2 Scorpii
Eta1 and Eta2 Coronae Australis
Winnecke 4 (Messier 40)
Uncertain
Albireo A and Albireo B
Kappa1 and Kappa2 Coronae Australis
Omicron1 and Omicron2 Centauri
Mizar system (Aa/Ab/Ba/Bb) and Alcor system (thus Mizar Ca/Cb), generally considered a physical system.
| Physical sciences | Stellar astronomy | Astronomy |
53620 | https://en.wikipedia.org/wiki/Puffin | Puffin | Puffins are any of three species of small alcids (auks) in the bird genus Fratercula. These are pelagic seabirds that feed primarily by diving in the water. They breed in large colonies on coastal cliffs or offshore islands, nesting in crevices among rocks or in burrows in the soil. Two species, the tufted puffin and horned puffin, are found in the North Pacific Ocean, while the Atlantic puffin is found in the North Atlantic Ocean.
All puffin species have predominantly black or black and white plumage, a stocky build, and large beaks that get brightly colored during the breeding season. They shed the colorful outer parts of their bills after the breeding season, leaving a smaller and duller beak. Their short wings are adapted for swimming with a flying technique underwater. In the air, they beat their wings rapidly (up to 400 times per minute) in swift flight, often flying low over the ocean's surface.
Etymology
The English name "puffin" – puffed in the sense of swollen – was originally applied to the fatty, salted meat of young birds of the unrelated Manx shearwater (Puffinus puffinus), formerly known as the "Manks puffin". Puffin is an Anglo-Norman word (Middle English pophyn or poffin) for the cured carcasses of nestling Manx shearwaters.
Taxonomy
The genus Fratercula was introduced by the French zoologist Mathurin Jacques Brisson in 1760 with the Atlantic puffin (Fratercula arctica) as the type species. The name Fratercula is Latin for "friar" from the word fraterculus "little brother", because the puffin's black and white plumage resemble robes worn by monks.
The genus contain three species. The rhinoceros auklet (Cerorhinca monocerata) has sometimes been included in the genus Fratercula, and some authors place the tufted puffin in the genus Lunda. The puffins and the rhinoceros auklet are closely related, together composing the subfamily Fraterculini.
The oldest alcid fossil is Hydrotherikornis from Oregon dating to the Late Eocene while fossils of Aethia and Uria go back to the Late Miocene. Molecular clocks have been used to suggest an origin in the Pacific in the Paleocene. Fossils from North Carolina were originally thought to have been of two Fratercula species, but were later reassigned to one Fratercula, the tufted puffin, and a Cerorhinca species. Another extinct species, Dow's puffin (Fratercula dowi) was found on the Channel Islands of California until the Late Pleistocene or early Holocene.
The Fraterculini are thought to have originated in the Pacific, primarily due to their greater diversity there; there exists only one extant species in the Atlantic, compared to two in the Pacific. This species has shown some significant signs of animal intelligence. In January 2020, some researchers reported that, Atlantic puffins were seen using sticks as a tool to scratch themselves. The Fraterculini fossil record in the Pacific extends at least as far back as the middle Miocene, with three fossil species of Cerorhinca, and material tentatively referred to that genus, in the middle Miocene to late Pliocene of southern California and northern Mexico. Although there are no records from the Miocene in the Atlantic, a re-examination of the North Carolina material indicated that the diversity of puffins in the early Pliocene was as great in the Atlantic as it is in the Pacific today. This diversity was achieved through influxes of puffins from the Pacific; the later loss of species was due to major oceanographic changes in the late Pliocene due to closure of the Panamanian Seaway and the onset of severe glacial cycles in the North Atlantic.
Extant species
Fossils
Description
The puffins are stocky, short-winged, and short-tailed birds, with black upper parts and white or brownish-grey underparts. The head has a black cap, the face is mainly white, and the feet are orange-red. The bill appears large and colorful during the breeding season. The colorful outer part of the bill is shed after the breeding season, revealing a smaller and duller true bill beneath. Because of their striking appearance they are also referred to as "clowns of the sea" and "sea parrots".
Although the puffins are vocal at their breeding colonies, they are silent at sea. They fly relatively high above the water, typically as compared with the of other auks.
Behaviour
Breeding
Puffins breed in colonies on coasts and islands; several current or former island breeding sites are referred to as Puffin Island. The male Atlantic puffin builds the nest and exhibits strong nest-site fidelity. Both sexes of the horned puffin help to construct their nest. Horned puffin burrows are usually about deep, ending in a chamber, while the tunnel leading to a tufted puffin burrow may be up to long. The nesting substrate of the tufted and Atlantic puffins is soft soil, into which tunnels are dug; in contrast, the nesting sites of horned puffins are rock crevices on cliffs. The Atlantic puffin burrow is usually lined with material such as grass, leaves, and feathers but is occasionally unlined. The eggs of the Atlantic puffin are typically creamy white but the occasional egg is tinged lilac.
Where rabbits breed, sometimes Atlantic puffins breed in rabbit burrows.
Puffins form long-term pair bonds or relationships. The female lays a single egg, and both parents incubate the egg and feed the chick (or "puffling"). The incubating parent holds the egg against its brood patch with its wings. The chicks fledge at night. After fledging, the chicks spend the first few years of their lives at sea, returning to breed about five years later. Puffins in captivity have been known to breed as early as three years of age.
After breeding, all three puffin species winter at sea, usually far from coasts and often extending south of the breeding range.
Iceland is the home to most of the Atlantic puffins with about 10 million individuals. The largest single puffin colony in the world is in the Westmann Isles of Iceland. In 2009, scientists estimated the number of nests to be 1.1 million, and number of individuals there is estimated to be up to 4 million.
Feeding
Like many auks, puffins eat both fish and zooplankton but feed their chicks primarily with small marine fish several times a day. The puffins are distinct in their ability to hold several (sometimes over a dozen) small fish at a time, crosswise in their bill, rather than regurgitating swallowed fish. This allows them to take longer foraging trips since they can come back with more food energy for their chick than a bird that can only carry one fish at a time. This behavior is made possible by the unique hinging mechanism of their beak, which allows the upper and lower biting edges to meet at any of a number of angles.
In 2019, animal experts observed puffins, in two separate geographic locations, using sticks to scratch themselves indicating that the seabirds have a basic ability to use tools.
Relationships with humans
Hunting
Puffins are hunted for eggs, feathers, and meat. Atlantic puffin populations drastically declined due to habitat destruction and exploitation during the 19th century and early 20th century. They continue to be hunted in Iceland and the Faroe Islands.
The Blasket Islands off the Irish coast of County Kerry saw a serious decline due to harvesting. Until the islands were abandoned in 1953, the islanders often lived just above starvation level. As a result, the puffins were hunted in large numbers for food.
The Atlantic puffin forms part of the national diet in Iceland, where the species does not have legal protection. Puffins are hunted by a technique called "sky fishing", which involves catching the puffins in a large net as they dive into the sea. Their meat is commonly featured on hotel menus. The fresh heart of a puffin is eaten raw as a traditional Icelandic delicacy. On the small Icelandic island of Grimsey as many as 200 puffins can be caught in a single morning.
Related places and products
The name of the English island Lundy is believed to come from the old Norse word for "puffin island" (Lundey), although an alternative explanation has been suggested with Lund referring to a copse, or wooded area.
The Atlantic puffin is the provincial bird of the Canadian province of Newfoundland and Labrador.
| Biology and health sciences | Charadriiformes | Animals |
53622 | https://en.wikipedia.org/wiki/Auk | Auk | Auks or alcids are birds of the family Alcidae in the order Charadriiformes. The alcid family includes the murres, guillemots, auklets, puffins, and murrelets. The family contains 25 extant or recently extinct species that are divided into 11 genera. Auks are found throughout the Northern Hemisphere.
Apart from the extinct great auk, all auks can fly, and are excellent swimmers and divers (appearing to "fly" in water), but their walking appears clumsy.
Names
Several species have different English names in Europe and North America. The two species known as "murres" in North America are called "guillemots" in Europe, and the species called little auk in Europe is referred to as dovekie in North America.
Etymology
The word "auk" is derived from Icelandic álka and Norwegian alka or alke from Old Norse ālka from Proto-Germanic *alkǭ (sea-bird, auk).
Taxonomy
The family name Alcidae comes from the genus Alca given by Carl Linnaeus in 1758 for the razorbill (Alca torda) from the Norwegian word alke.
Description
Auks are superficially similar to penguins, having black-and-white colours, upright posture, and some of their habits. Nevertheless, they are not closely related to penguins, but rather are believed to be an example of moderate convergent evolution. Auks are monomorphic (males and females are similar in appearance).
Extant auks range in size from the least auklet, at 85 g (3 oz) and , to the thick-billed murre, at and . Due to their short wings, auks have to flap their wings very quickly to fly.
Although not to the extent of penguins, auks have largely sacrificed flight, and also mobility on land, in exchange for swimming ability; their wings are a compromise between the best possible design for diving and the bare minimum needed for flying. This varies by subfamily, with the Uria guillemots (including the razorbill) and murrelets being the most efficient under the water, whereas the puffins and auklets are better adapted for flying and walking.
Feeding and ecology
The feeding behaviour of auks is often compared to that of penguins; both groups are wing-propelled, pursuit divers. In the region where auks live, their only seabird competition are cormorants (which are dive-powered by their strong feet). In areas where the two groups feed on the same prey, the auks tend to feed further offshore. Strong-swimming murres hunt faster, schooling fish, whereas auklets take slower-moving krill. Time depth recorders on auks have shown that they can dive as deep as in the case of Uria guillemots, for the Cepphus guillemots and for the auklets.
Breeding and colonies
Auks are pelagic birds, spending the majority of their adult lives on the open sea and going ashore only for breeding, although some species, such as the common guillemot, spend a great part of the year defending their nesting spot from others.
Auks are monogamous, and tend to form lifelong pairs. They typically lay a single egg, and they use the nesting site year after year.
Some species, such as the Uria guillemots (murres), nest in large colonies on cliff edges; others, such as the Cepphus guillemots, breed in small groups on rocky coasts; and the puffins, auklets, and some murrelets nest in burrows. All species except the Brachyramphus murrelets are colonial.
Evolution and distribution
Traditionally, the auks were believed to be one of the earliest distinct charadriiform lineages due to their characteristic morphology, but genetic analyses have demonstrated that these peculiarities are the product of strong natural selection, instead; as opposed to, for example, plovers (a much older charadriiform lineage), auks radically changed from a wading shorebird to a diving seabird lifestyle. Thus today, the auks are no longer separated in their own suborder (Alcae), but are considered part of the Lari suborder, which otherwise contains gulls and similar birds. Judging from genetic data, their closest living relatives appear to be the skuas, with these two lineages separating about 30 million years ago (Mya). Alternatively, auks may have split off far earlier from the rest of the Lari and undergone strong morphological, but slow genetic evolution, which would require a very high evolutionary pressure, coupled with a long lifespan and slow reproduction.
The earliest unequivocal fossils of auks are from the late Eocene, some 35 Mya. The genus Miocepphus, (from the Miocene, 15 Mya) is the earliest known from good specimens. Two very fragmentary fossils are often assigned to the Alcidae, although this may not be correct: Hydrotherikornis (Late Eocene) and Petralca (Late Oligocene). Most extant genera are known to exist since the Late Miocene or Early Pliocene (about 5 Mya). Miocene fossils have been found in both California and Maryland, but the greater diversity of fossils and tribes in the Pacific leads most scientists to conclude they first evolved there, and in the Miocene Pacific, the first fossils of extant genera are found. Early movement between the Pacific and the Atlantic probably happened to the south (since no northern opening to the Atlantic existed), with later movements across the Arctic Ocean. The flightless subfamily Mancallinae, which was apparently restricted to the Pacific Coast of southern North America and became extinct in the Early Pleistocene, is sometimes included in the family Alcidae under some definitions. One species, Miomancalla howardae, is the largest charadriiform of all time.
The family contains 25 extant or recently extinct species that are divided into 11 genera. The extant auks (subfamily Alcinae) are broken up into two main groups - the usually high-billed puffins (tribe Fraterculini) and auklets (tribe Aethiini), as opposed to the more slender-billed murres and true auks (tribe Alcini), and the murrelets and guillemots (tribes Brachyramphini and Cepphini). The tribal arrangement was originally based on analyses of morphology and ecology. mtDNA cytochrome b sequences, and allozyme studies confirm these findings except that the Synthliboramphus murrelets should be split into a distinct tribe, as they appear more closely related to the Alcini; in any case, assumption of a closer relationship between the former and the true guillemots was only weakly supported by earlier studies.
Of the genera, only a few species are placed in each. This is probably a product of the rather small geographic range of the family (the most limited of any seabird family), and the periods of glacial advance and retreat that have kept the populations on the move in a narrow band of subarctic ocean.
Today, as in the past, the auks are restricted to cooler northern waters. Their ability to spread further south is restricted as their prey hunting method, pursuit diving, becomes less efficient in warmer waters. The speed at which small fish (which along with krill are the auk's principal prey) can swim doubles as the temperature increases from , with no corresponding increase in speed for the bird. The southernmost auks, in California and Mexico, can survive there because of cold upwellings. The current paucity of auks in the Atlantic (six species), compared to the Pacific (19–20 species) is considered to be because of extinctions to the Atlantic auks; the fossil record shows many more species were in the Atlantic during the Pliocene. Auks also tend to be restricted to continental-shelf waters and breed on few oceanic islands.
Hydotherikornis oregonus (Described by Miller in 1931), the oldest purported alcid from the Eocene of California, is actually a petrel (as reviewed by Chandler in 1990) and is reassigned to the tubenoses (Procellariiformes). A 2003 paper, "The Earliest North American Record of Auk (Aves: Alcidae) From the Late Eocene of Central Georgia", reports a Late Eocene, wing-propelled, diving auk from the Priabonain stage of the Late Eocene. These sediments have been dated through Chandronian NALMA {North American Land Mammal Age}, at an estimate of 34.5 to 35.5 million years on the Eocene time scale for fossil-bearing sediments of the Clinchfield Formation, Gordon, Wilkinson County, Georgia. Furthermore, the sediments containing this unabraded portion of a left humerus (43.7 mm long) are tropical or subtropical as evidenced by a wealth of warm-water shark teeth, palaeophied snake vertebrae, and turtles.
Systematics
Basal and incertae sedis
Miocepphus (fossil: Middle Miocene of CE USA)
Miocepphus mcclungi Wetmore, 1940
Miocepphus bohaskai Wijnker and Olson, 2009
Miocepphus blowi Wijnker and Olson, 2009
Miocepphus mergulellus Wijnker and Olson, 2009
Subfamily Alcinae
Tribe Alcini – typical auks and murres
Uria
Common murre or common guillemot, Uria aalge
Thick-billed murre or, Brünnich's guillemot, Uria lomvia
Alle
Little auk or dovekie, Alle alle
Pinguinus
Great auk, Pinguinus impennis (extinct, c.1844)
Alca
Razorbill, Alca torda
Tribe Synthliboramphini – synthliboramphine murrelets
Synthliboramphus
Scripps's murrelet, Synthliboramphus scrippsi – formerly in S. hypoleucus ("Xantus's murrelet")
Guadalupe murrelet, Synthliboramphus hypoleucus – sometimes separated in Endomychura
Craveri's murrelet, Synthliboramphus craveri – sometimes separated in Endomychura
Ancient murrelet, Synthliboramphus antiquus
Japanese murrelet, Synthliboramphus wumizusume
Tribe Cepphini – true guillemots
Cepphus
Black guillemot or tystie, Cepphus grylle
Pigeon guillemot, Cepphus columba
Kurile guillemot, Cepphus columba snowi
Spectacled guillemot, Cepphus carbo
Tribe Brachyramphini – brachyramphine murrelets
Brachyramphus
Marbled murrelet, Brachyramphus marmoratus
Long-billed murrelet, Brachyramphus perdix
Kittlitz's murrelet, Brachyramphus brevirostris
Subfamily Fraterculinae
Tribe Aethiini – auklets
Ptychoramphus
Cassin's auklet, Ptychoramphus aleuticus
Aethia
Parakeet auklet, Aethia psittacula
Crested auklet, Aethia cristatella
Whiskered auklet, Aethia pygmaea
Least auklet, Aethia pusilla
Tribe Fraterculini – puffins
Cerorhinca
Rhinoceros auklet, Cerorhinca monocerata
Fratercula
Atlantic puffin, Fratercula arctica
Horned puffin, Fratercula corniculata
Tufted puffin, Fratercula cirrhata
Biodiversity of auks seems to have been markedly higher during the Pliocene. See the genus accounts for prehistoric species.
| Biology and health sciences | Charadriiformes | null |
53631 | https://en.wikipedia.org/wiki/Grylloblattidae | Grylloblattidae | Grylloblattidae, commonly known as the icebugs or ice crawlers, is a family of extremophile (psychrophile) and wingless insects that live in the cold on top of mountains and the edges of glaciers. It is the only member of Grylloblattodea, which is generally considered an order. Alternatively, Grylloblattodea, along with Mantophasmatodea (rock crawlers), have been ranked as suborders of the order Notoptera. Grylloblattids are wingless insects mostly less than 3 cm long, with a head resembling that of a cockroach, with long antennae and having elongated cerci arising from the tip of their abdomen. They cannot tolerate warmth (most species will die at 10 °C) and many species have small distribution ranges.
Overview
Grylloblattids, ice crawlers or icebugs puzzled the scientists who discovered them in 1914, E.M. Walker and T.B. Kurata; the first species named was Grylloblatta campodeiformis, which means "cricket-cockroach shaped like a Campodea" (a kind of two-pronged bristletail). Most are nocturnal and appear to feed on detritus. They have long antennae (23–45 segments) and long cerci (5–8 segments), but no wings. Their eyes are either missing or reduced and they have no ocelli (simple eyes). Their closest living relatives are the recently discovered Mantophasmatodea. Most species are less than 3 cm long, the largest being Namkungia magnus.
The family has its own order, Grylloblattodea (sometimes considered a suborder of Notoptera). It contains 5 genera and about 34 extant species.
Most species have restricted distributions and small populations and with increased warming their habitats are threatened, making them endangered. In North America some species like Grylloblatta barberi and G. oregonensis are known from single sites.
Habitat and distribution
Grylloblattodea are nocturnal extremophiles typically found in leaf litter and under stones in extremely cold environments, usually at higher elevations. They are known to inhabit cold temperate forests to glaciers and the edges of ice sheets. Their optimal living temperature is between . They can be killed at colder temperatures due to ice formation in the body, so when the temperature drops below their optimal range they survive by living under snow pack near the soil. They have a very narrow range of temperatures that they prefer and cannot withstand high temperatures; many species are killed when the temperature rises about 5 °C above their optimal temperature. They move in response to the seasons so as to maintain an optimal temperature in their foraging habitat.
Grylloblattidae are patchily distributed in glaciers, caves, montane environments, and occasionally lower-elevation forests in western North America, East Asia (Korea and Japan), and Central Asia (Siberia, China, and Kazakhstan). They are predicted to occur in several other mountain chains in Asia, including parts of the Himalayas.
Diet
They are omnivorous, but feed primarily on dead arthropods and carrion. When arthropod carcasses are scarce, they subsist on plant material.
Evolution
Grylloblattidae is generally thought to have emerged from within the "Grylloblattida", a poorly defined group of extinct winged insects that first appeared in the Late Carboniferous, over 300 million years ago. The winged Aristovia from the mid-Cretaceous Burmese amber of Myanmar, around 100 million years ago is thought to be closely related to modern Grylloblattidae due to its very similar mouthparts.
Taxonomy
List of Grylloblattodea genera and species along with their type localities:
Galloisiana Caudell 1924 – Far East Asia
Galloisiana chujoi Gurney 1961 – type locality: Oninoiwaya Cave, Japan
Galloisiana kiyosawai Asahina 1959 – type locality: Hirayu-Onsen, Japan
Galloisiana kosuensis Namkung 1974 – type locality: Gosu Cave, South Korea
Galloisiana nipponensis (Caudell & King 1924) – type locality: Lake Chūzenji, Japan
Galloisiana notabilis (Silvestri 1927) – type locality: Nagasaki Prefecture, Japan
Galloisiana odaesanensis Kim & Lee 2007 – type locality: Mount Odae, South Korea
Galloisiana olgae Vrsansky & Storozhenko 2001 – type locality: Mount Olga, Russia
Galloisiana sinensis Wang 1987 – type locality: Changbaishan, Jilin, PR China
Galloisiana sofiae Szeptycki 1987 – type locality: Mount Myoyang, South Korea
Galloisiana ussuriensis Storozhenko 1988 – type locality: Primorsky Krai, Russia
Galloisiana yezoensis Asahina 1961 – type locality: Miyazaki-Toge, Japan
Galloisiana yuasai Asahina 1959 – type locality: Tokugo-Toge, Japan
Grylloblatta Walker 1914 – western North America
Grylloblatta barberi Caudell 1924 – type locality: Sunny Side Mine, Mount Lassen area, California, USA
Grylloblatta bifratrilecta Gurney 1953 – type locality: Sonora Pass, California, USA
Grylloblatta campodeiformis Walker 1914 – type locality: Sulphur Mountain, Alberta, Canada
Grylloblatta chandleri Kamp 1963 – type locality: Eagle Lake, California, USA
Grylloblatta chintimini Marshall & Lytle 2015 – type locality: Marys Peak, Oregon, USA
Grylloblatta chirurgica Gurney 1961 – type locality: Ape Cave, Washington, USA
Grylloblatta gurneyi Kamp 1963 – type locality: Lava Beds National Monument, California, USA
Grylloblatta marmoreus Schoville 2012 – type locality: Marble Mountains, California, USA
Grylloblatta newberryensis Marshall & Lytle 2015 – type locality: Newberry Volcano, Oregon, USA
Grylloblatta oregonensis Schoville 2012 – type locality: Oregon Caves National Monument, USA
Grylloblatta rothi Gurney 1953 – type locality: Happy Valley, Deschutes County, Oregon, USA
Grylloblatta scudderi Kamp 1979 – type locality: Mount Paul, Alberta, Canada
Grylloblatta sculleni Gurney 1937 – type locality: Scott Camp, Deschutes County, Oregon, USA
Grylloblatta siskiyouensis Schoville 2012 – type locality: Oregon Caves National Monument, USA
Grylloblatta washoa Gurney 1961 – type locality: Echo Summit, California, USA
Grylloblattella Storozhenko 1988 – Central Asia
Grylloblattella cheni Bai, Wang & Yang 2010 – type locality: Ake Kule Lake, Xinjiang, China
Grylloblattella pravdini (Storozhenko & Oliger 1984) – type locality: Teletskoye Lake, Russia
Grylloblattella sayanensis Storozhenko 1996 – type locality: Sambyl Pass, Russia
Grylloblattina Bey-Bienko 1951 – Far East Asia
Grylloblattina djakonovi Bey-Bienko 1951
Namkungia Storozhenko & Park 2002 – Korea
Namkungia biryongensis (Namkung 1974) – type locality: Biryong Cave, South Korea
Namkungia magna (Namkung 1986) – type locality: Balgudeok Cave, South Korea
In total, there are 33 extant species and 5 extant genera described as of 2015.
| Biology and health sciences | Insects: General | Animals |
53635 | https://en.wikipedia.org/wiki/Odonata | Odonata | Odonata is an order of predatory flying insects that includes the dragonflies and damselflies (as well as the Epiophlebia damsel-dragonflies). The two major groups are distinguished with dragonflies (Anisoptera) usually being bulkier with large compound eyes together and wings spread up or out at rest, while damselflies (suborder Zygoptera) are usually more slender with eyes placed apart and wings folded together along body at rest. Adult odonates can land and perch, but rarely walk.
All odonates have aquatic larvae called naiads or nymphs, and all of them, larvae and adults, are carnivorous and are almost entirely insectivorous, although at the larval stage they will eat anything that they can overpower, including small fish, tadpoles, and even adult newts. The adults are superb aerial hunters and their legs are specialised for catching prey in flight.
Odonata in its narrow sense forms a subgroup of the broader Odonatoptera, which contains other dragonfly-like insects.
The scientific study of the Odonata is called odonatology.
Etymology and terminology
Johan Christian Fabricius coined the term Odonata in 1793 from the Ancient Greek (Ionic form of ) "tooth". One hypothesis is that it was because their maxillae are notably toothed.
The word dragonfly usually denotes only Anisoptera, but is sometimes used to mean all odonatans. Odonata enthusiasts avoid ambiguity by using the term true dragonfly, or simply anisopteran, when they mean just the Anisoptera. An alternative term warriorfly has been proposed.
External morphology
Size
The largest living odonate is the giant Central American helicopter damselfly Megaloprepus coerulatus (Zygoptera: Pseudostigmatidae) with a wing span of . The heaviest living odonates are Tetracanthagyna plagiata (Anisoptera: Aeshnidae) with a wing span of , and Petalura ingentissima (Anisoptera: Petaluridae) with a body length of (some sources ) and wing span of . The longest extant odonate is the Neotropical helicopter damselfly Mecistogaster linearis (Zygoptera: Pseudostigmatidae) with a body length of .
The smallest living dragonfly is Nannophya pygmaea (Anisoptera: Libellulidae) from east Asia, with a body length of and a wing span of . The smallest damselflies (and also the smallest odonates) are species of the genus Agriocnemis (Zygoptera: Coenagrionidae) with a wing span of only .
Description
These insects characteristically have large rounded heads covered mostly by well-developed, compound eyes, which provide good vision, legs that facilitate catching prey (other insects) in flight, two pairs of long, transparent wings that move independently, and elongated abdomens. They have three ocelli and short antennae. The mouthparts are on the underside of the head and include simple chewing mandibles in the adult.
Flight in the Odonata is direct, with flight muscles attaching directly to the wings; rather than indirect, with flight muscles attaching to the thorax, as is found in the Neoptera. This allows active control of the amplitude, frequency, angle of attack, camber and twist of each of the four wings entirely independently.
In most families, there is a structure on the leading edge near the tip of the wing called the pterostigma. This is a thickened, hemolymph–filled and often colorful area bounded by veins. The functions of the pterostigma are not fully known, but it most probably has an aerodynamic effect and may also have a visual function. More mass at the end of the wing may also reduce the energy needed to move the wings up and down. The right combination of wing stiffness and wing mass could reduce the energy consumption of flying. A pterostigma is also found among other insects, such as bees.
The nymphs have stockier, shorter, bodies than the adults. In addition to lacking wings, their eyes are smaller, their antennae longer, and their heads are less mobile than in the adult. Their mouthparts are modified, with the labium being adapted into a unique prehensile organ called a labial mask for grasping prey. Damselfly nymphs breathe through external gills on the abdomen, while dragonfly nymphs respire through an organ in their rectum.
Evolution
Fossil history
Members of the crown group Odonata first appeared during the Late Triassic, though members of their total group, Odonatoptera, first appeared in the Late Carboniferous, making them one of the earliest groups of winged insects. The fossils of odonates and their cousins, including Paleozoic "giant dragonflies" like Meganeuropsis permiana from the Permian of North America, reached wing spans of up to and a body length of , making it the largest insect of all time. This insect belonged to the order Meganisoptera, the griffinflies, related to odonates but not part of the modern order Odonata in the restricted sense. They have one of the most complete fossil records going back 319 million years.
The Odonata is closely related to mayflies and several extinct orders in a group called the Palaeoptera, but this grouping might be paraphyletic. What they do share with mayflies is the nature of how the wings are articulated and held in rest.
Tarsophlebiidae is a prehistoric family of Odonatoptera that can be considered either a basal lineage of Odonata or their immediate sister taxon.
Phylogeny
The phylogenetic tree of the orders and suborders of odonates according to Bybee et al. 2021:
Taxonomy
In some treatments, the Odonata are understood in an expanded sense, essentially synonymous with the superorder Odonatoptera, but not including the prehistoric Protodonata. In this approach, instead of Odonatoptera, the term Odonatoidea is used. The systematics of the "Palaeoptera" are by no means resolved; what can be said however is that regardless of whether they are called "Odonatoidea" or "Odonatoptera", the Odonata and their extinct relatives do form a clade.
The Anisoptera was long treated as a suborder, with a third suborder, the Anisozygoptera (ancient dragonflies). However, the combined suborder Epiprocta (in which Anisoptera is an infraorder) was proposed when it was thought that the "Anisozygoptera" was paraphyletic, composed of mostly extinct offshoots of dragonfly evolution. The four living species placed in that group are (in this treatment) in the infraorder Epiophlebioptera, whereas the fossil taxa that were formerly there are now dispersed about the Odonatoptera (or Odonata sensu lato). World Odonata List considers Anisoptera as a suborder along with Zygoptera and Anisozygoptera as well-understood and widely preferred terms.
Cladogram of Epiprocta after Rehn et al. 2003:Cladogram of Odonatoptera including Odonata by Deregnaucourt et al. 2023.
Ecology and life cycle
Odonates are aquatic or semi-aquatic as juveniles. Thus, adults are most often seen near bodies of water and are frequently described as aquatic insects. However, many species range far from water. They are carnivorous (or more specifically insectivorous) throughout their life, mostly feeding on smaller insects.
Male Odonata have complex genitalia, different from those found in other insects. These include grasping cerci at the tip of the abdomen for holding the female, and a secondary set of copulatory organs located between the second and third abdominal segment in which the spermatozoa are stored after being produced by the primary genitals— whose external opening is known as the genital pore, on the ninth abdominal segment. This process is called intra-male sperm translocation (ST). Because the male copulatory organ has evolved independently from that in other insects, it has been suggested the stem-group dragonflies had external sperm transfer. To mate, the male claspers grasps the female by the thorax (Zygoptera) or head (Anisoptera) while the female bends her abdomen so that her own genitalia can be grasped by the copulatory organ holding the sperm. This is known as the "wheel" position. In Anisoptera, males often mate while flying, lifting the females in the air, which typically last from a couple of seconds to a minute or two, whereas the males in Zygoptera mate while perched. They might even move to different spots during the mating process, which can make it last longer, anywhere between five and ten minutes. Male Odonata are very competitive when it comes to mating that in some species, the males use the cerci located at the tip of the abdomen to remove the sperm of a rival male's from the female and put in his own.
Eggs are laid in water or on vegetation near water or wet places, and hatch to produce pronymphs which live off the nutrients that were in the egg. They then develop into instars with approximately 9–14 molts that are (in most species) voracious predators on other aquatic organisms, including small fishes. The nymphs grow and molt, usually in dusk or dawn, into the flying teneral immature adults, whose color is not yet developed. These transform into reproductive adults.
Odonates can act as bioindicators of water quality in rivers because they rely on high quality water for proper development in early life. Since their diet consists entirely of insects, odonate density is directly proportional to the population of prey, and their abundance indicates the abundance of prey in the examined ecosystem. Species richness of vascular plants has also been positively correlated with the species richness of dragonflies in a given habitat. This means that in a location such as a lake, if one finds a wide variety of odonates, then a similarly wide variety of plants should also be present. This correlation is not common to all bioindicators, as some may act as indicators for a different environmental factor, such as the pool frog acting as a bioindicator of water quality due to its high quantity of time spent in and around water.
In addition, odonates are very sensitive to changes to average temperature. Many species have moved to higher elevations and latitudes as global temperature rises and habitats dry out. Changes to the life cycle have been recorded with increased development of the instar stages and smaller adult body size as the average temperature increases. As the territory of many species starts to overlap, the rate hybridization of species that normally do not come in contact is increasing. If global climate change continues many members of Odonata will start to disappear. Because odonates are such an old order and have such a complete fossil record they are an ideal species to study insect evolution and adaptation. For example, they are one of the first insects to develop flight and it is likely that this trait only evolved once in insects, looking at how flight works in odonates, the rest of flight can be mapped out.
Cannibalism
Cannibalism has been recorded in many species of odonates, at both the larval and adult stages. Cannibalism is caused by either errors in species recognition, intrasexual competition for mating, or prevention of mating harassment.
Gallery
| Biology and health sciences | Odonata | null |
53662 | https://en.wikipedia.org/wiki/Peccary | Peccary | Peccaries (also javelinas or skunk pigs) are pig-like ungulates of the family Tayassuidae (New World pigs). They are found throughout Central and South America, Trinidad in the Caribbean, and in the southwestern area of North America. Peccaries usually measure between in length, and a full-grown adult usually weighs about . They represent the closest relatives of the family Suidae, which contains pigs and relatives. Together Tayassuidae and Suidae are grouped in the suborder Suina within the order Artiodactyla (even-toed ungulates).
Peccaries are social creatures that live in herds. They are omnivores and eat roots, grubs, and a variety of other foods. They can identify each other by their strong odors. A group of peccaries that travel and live together is called a squadron. A squadron of peccaries averages between six and nine members.
Peccaries first appeared in North America during the Miocene and migrated into South America during the Pliocene–Pleistocene as part of the Great American Interchange.
They are often confused with feral domestic pigs, commonly known as "razorback" hogs in many parts of the United States, when the two occur in the wild in similar ranges.
The Maya kept herds of peccaries, using them in rituals and for food. They are kept as pets in many countries in addition to being raised on farms as a source of food.
Etymology
The word peccary is derived from the Carib word or .
In Portuguese, a peccary is called , , , , among other names like or . In Spanish, it is called , (a word also used to describe wild boar), , or . The word javelina derives from the Spanish word for "wild boar". In French Guiana and Suriname, the animal is called pakira.
The scientific name Tayassuidae derives from the same source as the Portuguese .
Characteristics
A peccary is a medium-sized animal, with a strong resemblance to a pig. Like a pig, it has a snout ending in a cartilaginous disc and eyes that are small relative to its head. Also like a pig, it uses only the middle two digits for walking, although, unlike pigs, the other toes may be altogether absent. Its stomach is not ruminating. Though it has three chambers, it is more complex than those of pigs. Peccaries are foregut fermenters (pigs are hindgut fermenters). This foregut fermentation, similar to but separately evolved from a ruminant, is an example of convergent evolution.
Peccaries are omnivores and will eat insects, grubs, and occasionally small animals, although their preferred foods consist of roots, grasses, seeds, fruit, and cacti—particularly prickly pear. Pigs and peccaries can be differentiated by a number of characteristic, including tails and ear shape. The ears of pigs are large and upright and often pointed while the ears of peccaries are small and rounded. Pigs also have tasseled tails, but peccaries' tails are small and discreet.
The most noticeable difference between pigs and peccaries is the shape of the canine teeth, or tusks. In European pigs, the tusks are long and curve around on themselves, whereas in peccaries, the tusks are short and straight and interlock with each other, prohibiting side-to-side movement of the jaw. The jaws and tusks of peccaries are adapted for crushing hard seeds and slicing into plant roots, and they also use their tusks to defend against predators. The dental formula for peccaries is:
By rubbing the tusks together, they can make a chattering noise that warns potential predators to stay away.
Peccaries are social animals, often forming herds. Over 100 individuals have been recorded for a single herd of white-lipped peccaries, but collared and Chacoan peccaries usually form smaller groups. Such social behavior seems to have been the situation in extinct peccaries as well. The giant peccary (Pecari maximus) of Brazil appears to be less social, primarily living in pairs. Peccaries rely on their social structure to defend territory, protect against predators, regulate temperature, and interact with other members of the species.
Peccaries have scent glands below each eye and another on their backs, though these are believed to be rudimentary in P. maximus. They use the scent to mark herd territories, which range from . They also mark other herd members with these scent glands by rubbing one against another. The pungent odor allows peccaries to recognize other members of their herd, despite their myopic vision. The odor is strong enough to be detected by humans, which earns the peccary the nickname of "skunk pig".
Species
Extant species
Three (possibly four) living species of peccaries are found from the Southwestern United States through Central America and into South America and Trinidad, each in their own genus.
Tayassu
White-lipped peccary (T. pecari)
Catagonus
Chacoan peccary (C. wagneri)
Dicotyles
Collared peccary (D. tajacu)
The collared peccary (Dicotyles tajacu) or "musk hog", referring to the animal's scent glands, occurs from the Southwestern United States into South America and the island of Trinidad. The coat consists of wiry peppered black, gray, and brown hair with a lighter colored "collar" circling the shoulders. They bear young year-round, but most often between November and March, with the average litter size consisting of two to three offspring. They are found in many habitats, from arid scrublands to humid tropical rain forests. The collared peccary is well-adapted to habitat disturbed by humans, merely requiring sufficient cover. They can be found in cities and agricultural land throughout their range.
Notable populations exist in the suburbs of Phoenix and Tucson, Arizona, where they feed on ornamental plants and other cultivated vegetation. There are also urban populations as far north as Sedona, Arizona, where they have been known to fill a niche similar to raccoons and other urban scavengers. In Arizona they are often called by their Spanish name "javelinas". Collared peccaries are generally found in bands of 8 to 15 animals of various ages. They defend themselves if they feel threatened, but otherwise tend to ignore humans.
A second species, the white-lipped peccary (Tayassu pecari), is mainly found in rainforests of Central and South America, but also known from a wide range of other habitats such as dry forests, grasslands, mangrove, cerrado, and dry xerophytic areas. The two main threats to their survival are deforestation and hunting.
The third species, the Chacoan peccary (Catagonus wagneri). It is found in the dry shrub habitat or Chaco of Paraguay, Bolivia, and Argentina. The Chacoan peccary has the distinction of having been first described based on fossils and was originally thought to be an extinct species. In 1975, the animal was discovered in the Chaco region of Paraguay. The species was well known to the native people.
A fourth as yet unconfirmed species, the giant peccary (Dicotyles maximus), was described from the Brazilian Amazon and north Bolivia by Dutch biologist Marc van Roosmalen. Though relatively recently discovered, it has been known to the local Tupi people as caitetu munde, which means "great peccary which lives in pairs". Thought to be the largest extant peccary, it can grow to in length. Its pelage is completely dark gray, with no collars whatsoever. Unlike other peccaries, it lives in pairs, or with one or two offspring. However, the scientific evidence for considering it as a species separate from the collared peccary has later been questioned, leading the IUCN to treat it as a synonym.
During the Late Pleistocene, two extinct peccaries, Mylohyus and Platygonus, were widespread across North America (and in the case of Platygonus, South America), but became extinct at the end of the Pleistocene around 12,000 years ago following the arrival of humans.
Extinct genera
In addition, Tayassuidae have a well-attested fossil record, and numerous extinct genera are known:
†Aptenohyus
†Cynorca
†Egatochoerus
†Floridachoerus
†Macrogenis
†Mckennahyus
†Mylohyus
†Platygonus
†Prochoerus
†Prosthennops
†Simojovelhyus
†Skinnerhyus
†Thinohyus
†Woodburnehyus
Evolution
Although some taxa from the Old World like the European Miocene Taucanamo have been suggested to be members of Tayussidae, their assignation to the group is equivocal, with a 2017 phylogenetic analysis recovering Taucanamo outside the clade containing suids and peccaries. The oldest unambiguous fossils of peccaries are from the Early Miocene of North America, with the North American Eocene-Oligocene genus Perchoerus, also often considered an early peccary, recovered outside the clade containing peccaries and suids.
Although common in South America today, peccaries did not reach there until about three million years ago during the Great American Interchange, when the Isthmus of Panama formed, connecting North America and South America. At that time, many North American animals—including peccaries, llamas and tapirs—entered South America, while some South American species, such as the ground sloths and opossums, migrated north. Several species of peccary across the genera Platygonus and Mylohyus remained in North America until their extinction following the colonization of the continent by humans via Beringia at the end of the Pleistocene. Today, 2 of the 3 species are relegated to the Neotropical realm, but the collared peccary ranges into northern Mexico and the southwestern United States.
Domestication
Peccaries bear a familial resemblance to true pigs due to their common ancestry, and are in the same suborder as swine (Suina). They have been present in South America since prehistoric times. The earliest scientific description of peccaries in the New World is in Brazil in 1547 and referred to them as "wild pigs".
It has been documented that peccaries were tamed, penned, and raised for food and ritual purposes in the Yucatán, Panama, the southern Caribbean, and Colombia at the time of the Conquest. Archaeological remains of peccaries have been found in Mesoamerica from the Preclassic (or Formative) period up until immediately before Spanish contact. Specifically, peccary remains have been found at Early Formative Olmec civilization sites.
The peccary is not readily suitable for modern captive breeding, lacking suitable characteristics for intensive or semi-intensive systems. Peccaries require a higher age before they are able to give birth (parturition) and have a tendency towards infanticide.
Relation with feral pigs
Recently established Brazilian boar populations are not to be confused with long-established populations of feral domestic pigs, which have existed mainly in the Pantanal for more than 100 years, along with native peccaries. The demographic dynamics of the interaction between feral pig populations and those of the two native species of peccaries (collared peccary and white-lipped peccary) is obscure and is still being studied. The existence of feral pigs could somewhat ease jaguar predation on peccary populations, as jaguars show a preference for hunting pigs when they are available.
| Biology and health sciences | Artiodactyla | null |
53664 | https://en.wikipedia.org/wiki/Kingdom%20%28biology%29 | Kingdom (biology) | In biology, a kingdom is the second highest taxonomic rank, just below domain. Kingdoms are divided into smaller groups called phyla (singular phylum).
Traditionally, textbooks from Canada and the United States have used a system of six kingdoms (Animalia, Plantae, Fungi, Protista, Archaea/Archaebacteria, and Bacteria or Eubacteria), while textbooks in other parts of the world, such as Bangladesh, Brazil, Greece, India, Pakistan, Spain, and the United Kingdom have used five kingdoms (Animalia, Plantae, Fungi, Protista and Monera).
Some recent classifications based on modern cladistics have explicitly abandoned the term kingdom, noting that some traditional kingdoms are not monophyletic, meaning that they do not consist of all the descendants of a common ancestor. The terms flora (for plants), fauna (for animals), and, in the 21st century, funga (for fungi) are also used for life present in a particular region or time.
Definition and associated terms
When Carl Linnaeus introduced the rank-based system of nomenclature into biology in 1735, the highest rank was given the name "kingdom" and was followed by four other main or principal ranks: class, order, genus and species. Later two further main ranks were introduced, making the sequence kingdom, phylum or division, class, order, family, genus and species. In 1990, the rank of domain was introduced above kingdom.
Prefixes can be added so subkingdom (subregnum) and infrakingdom (also known as infraregnum) are the two ranks immediately below kingdom. Superkingdom may be considered as an equivalent of domain or empire or as an independent rank between kingdom and domain or subdomain. In some classification systems the additional rank branch (Latin: ramus) can be inserted between subkingdom and infrakingdom, e.g., Protostomia and Deuterostomia in the classification of Cavalier-Smith.
History
Two kingdoms of life
The classification of living things into animals and plants is an ancient one. Aristotle (384–322 BC) classified animal species in his History of Animals, while his pupil Theophrastus (–) wrote a parallel work, the Historia Plantarum, on plants.
Carl Linnaeus (1707–1778) laid the foundations for modern biological nomenclature, now regulated by the Nomenclature Codes, in 1735. He distinguished two kingdoms of living things: Regnum Animale ('animal kingdom') and Regnum Vegetabile ('vegetable kingdom', for plants). Linnaeus also included minerals in his classification system, placing them in a third kingdom, Regnum Lapideum.
Three kingdoms of life
In 1674, Antonie van Leeuwenhoek, often called the "father of microscopy", sent the Royal Society of London a copy of his first observations of microscopic single-celled organisms. Until then, the existence of such microscopic organisms was entirely unknown. Despite this, Linnaeus did not include any microscopic creatures in his original taxonomy.
At first, microscopic organisms were classified within the animal and plant kingdoms. However, by the mid–19th century, it had become clear to many that "the existing dichotomy of the plant and animal kingdoms [had become] rapidly blurred at its boundaries and outmoded".
In 1860 John Hogg proposed the Protoctista, a third kingdom of life composed of "all the lower creatures, or the primary organic beings"; he retained Regnum Lapideum as a fourth kingdom of minerals. In 1866, Ernst Haeckel also proposed a third kingdom of life, the Protista, for "neutral organisms" or "the kingdom of primitive forms", which were neither animal nor plant; he did not include the Regnum Lapideum in his scheme. Haeckel revised the content of this kingdom a number of times before settling on a division based on whether organisms were unicellular (Protista) or multicellular (animals and plants).
Four kingdoms
The development of microscopy revealed important distinctions between those organisms whose cells do not have a distinct nucleus (prokaryotes) and organisms whose cells do have a distinct nucleus (eukaryotes). In 1937 Édouard Chatton introduced the terms "prokaryote" and "eukaryote" to differentiate these organisms.
In 1938, Herbert F. Copeland proposed a four-kingdom classification by creating the novel Kingdom Monera of prokaryotic organisms; as a revised phylum Monera of the Protista, it included organisms now classified as Bacteria and Archaea. Ernst Haeckel, in his 1904 book The Wonders of Life, had placed the blue-green algae (or Phycochromacea) in Monera; this would gradually gain acceptance, and the blue-green algae would become classified as bacteria in the phylum Cyanobacteria.
In the 1960s, Roger Stanier and C. B. van Niel promoted and popularized Édouard Chatton's earlier work, particularly in their paper of 1962, "The Concept of a Bacterium"; this created, for the first time, a rank above kingdom—a superkingdom or empire—with the two-empire system of prokaryotes and eukaryotes. The two-empire system would later be expanded to the three-domain system of Archaea, Bacteria, and Eukaryota.
Five kingdoms
The differences between fungi and other organisms regarded as plants had long been recognised by some; Haeckel had moved the fungi out of Plantae into Protista after his original classification, but was largely ignored in this separation by scientists of his time. Robert Whittaker recognized an additional kingdom for the Fungi. The resulting five-kingdom system, proposed in 1969 by Whittaker, has become a popular standard and with some refinement is still used in many works and forms the basis for new multi-kingdom systems. It is based mainly upon differences in nutrition; his Plantae were mostly multicellular autotrophs, his Animalia multicellular heterotrophs, and his Fungi multicellular saprotrophs.
The remaining two kingdoms, Protista and Monera, included unicellular and simple cellular colonies. The five kingdom system may be combined with the two empire system. In the Whittaker system, Plantae included some algae. In other systems, such as Lynn Margulis's system of five kingdoms, the plants included just the land plants (Embryophyta), and Protoctista has a broader definition.
Following publication of Whittaker's system, the five-kingdom model began to be commonly used in high school biology textbooks. But despite the development from two kingdoms to five among most scientists, some authors as late as 1975 continued to employ a traditional two-kingdom system of animals and plants, dividing the plant kingdom into subkingdoms Prokaryota (bacteria and cyanobacteria), Mycota (fungi and supposed relatives), and Chlorota (algae and land plants).
Six kingdoms
In 1977, Carl Woese and colleagues proposed the fundamental subdivision of the prokaryotes into the Eubacteria (later called the Bacteria) and Archaebacteria (later called the Archaea), based on ribosomal RNA structure; this would later lead to the proposal of three "domains" of life, of Bacteria, Archaea, and Eukaryota. Combined with the five-kingdom model, this created a six-kingdom model, where the kingdom Monera is replaced by the kingdoms Bacteria and Archaea. This six-kingdom model is commonly used in recent US high school biology textbooks, but has received criticism for compromising the current scientific consensus. But the division of prokaryotes into two kingdoms remains in use with the recent seven kingdoms scheme of Thomas Cavalier-Smith, although it primarily differs in that Protista is replaced by Protozoa and Chromista.
Eight kingdoms
Thomas Cavalier-Smith supported the consensus at that time, that the difference between Eubacteria and Archaebacteria was so great (particularly considering the genetic distance of ribosomal genes) that the prokaryotes needed to be separated into two different kingdoms. He then divided Eubacteria into two subkingdoms: Negibacteria (Gram-negative bacteria) and Posibacteria (Gram-positive bacteria). Technological advances in electron microscopy allowed the separation of the Chromista from the Plantae kingdom. Indeed, the chloroplast of the chromists is located in the lumen of the endoplasmic reticulum instead of in the cytosol. Moreover, only chromists contain chlorophyll c. Since then, many non-photosynthetic phyla of protists, thought to have secondarily lost their chloroplasts, were integrated into the kingdom Chromista.
Finally, some protists lacking mitochondria were discovered. As mitochondria were known to be the result of the endosymbiosis of a proteobacterium, it was thought that these amitochondriate eukaryotes were primitively so, marking an important step in eukaryogenesis. As a result, these amitochondriate protists were separated from the protist kingdom, giving rise to the, at the same time, superkingdom and kingdom Archezoa. This superkingdom was opposed to the Metakaryota superkingdom, grouping together the five other eukaryotic kingdoms (Animalia, Protozoa, Fungi, Plantae and Chromista). This was known as the Archezoa hypothesis, which has since been abandoned; later schemes did not include the Archezoa–Metakaryota divide.
‡ No longer recognized by taxonomists.
Six kingdoms (1998)
In 1998, Cavalier-Smith published a six-kingdom model, which has been revised in subsequent papers. The version published in 2009 is shown below. Cavalier-Smith no longer accepted the importance of the fundamental Eubacteria–Archaebacteria divide put forward by Woese and others and supported by recent research. The kingdom Bacteria (sole kingdom of empire Prokaryota) was subdivided into two sub-kingdoms according to their membrane topologies: Unibacteria and Negibacteria. Unibacteria was divided into phyla Archaebacteria and Posibacteria; the bimembranous-unimembranous transition was thought to be far more fundamental than the long branch of genetic distance of Archaebacteria, viewed as having no particular biological significance.
Cavalier-Smith does not accept the requirement for taxa to be monophyletic ("holophyletic" in his terminology) to be valid. He defines Prokaryota, Bacteria, Negibacteria, Unibacteria, and Posibacteria as valid paraphyla (therefore "monophyletic" in the sense he uses this term) taxa, marking important innovations of biological significance (in regard of the concept of biological niche).
In the same way, his paraphyletic kingdom Protozoa includes the ancestors of Animalia, Fungi, Plantae, and Chromista. The advances of phylogenetic studies allowed Cavalier-Smith to realize that all the phyla thought to be archezoans (i.e. primitively amitochondriate eukaryotes) had in fact secondarily lost their mitochondria, typically by transforming them into new organelles: Hydrogenosomes. This means that all living eukaryotes are in fact metakaryotes, according to the significance of the term given by Cavalier-Smith. Some of the members of the defunct kingdom Archezoa, like the phylum Microsporidia, were reclassified into kingdom Fungi. Others were reclassified in kingdom Protozoa, like Metamonada which is now part of infrakingdom Excavata.
Because Cavalier-Smith allows paraphyly, the diagram below is an 'organization chart', not an 'ancestor chart', and does not represent an evolutionary tree.
Seven kingdoms
Cavalier-Smith and his collaborators revised their classification in 2015. In this scheme they introduced two superkingdoms of Prokaryota and Eukaryota and seven kingdoms. Prokaryota have two kingdoms: Bacteria and Archaea. (This was based on the consensus in the Taxonomic Outline of Bacteria and Archaea, and the Catalogue of Life). The Eukaryota have five kingdoms: Protozoa, Chromista, Plantae, Fungi, and Animalia. In this classification a protist is any of the eukaryotic unicellular organisms.
Summary
The kingdom-level classification of life is still widely employed as a useful way of grouping organisms, notwithstanding some problems with this approach:
Kingdoms such as Protozoa represent grades rather than clades, and so are rejected by phylogenetic classification systems.
The most recent research does not support the classification of the eukaryotes into any of the standard systems. In 2009, Andrew Roger and Alastair Simpson emphasized the need for diligence in analyzing new discoveries: "With the current pace of change in our understanding of the eukaryote tree of life, we should proceed with caution." Kingdoms are rarely used in academic phylogeny and are more common in introductory education, where 5-6 kingdom models are preferred.
Beyond traditional kingdoms
While the concept of kingdoms continues to be used by some taxonomists, there has been a movement away from traditional kingdoms, as they are no longer seen as providing a cladistic classification, where there is emphasis in arranging organisms into natural groups.
Three domains of life
Based on RNA studies, Carl Woese thought life could be divided into three large divisions and referred to them as the "three primary kingdom" model or "urkingdom" model.
In 1990, the name "domain" was proposed for the highest rank. This term represents a synonym for the category of dominion (lat. dominium), introduced by Moore in 1974. Unlike Moore, Woese et al. (1990) did not suggest a Latin term for this category, which represents a further argument supporting the accurately introduced term dominion.
Woese divided the prokaryotes (previously classified as the Kingdom Monera) into two groups, called Eubacteria and Archaebacteria, stressing that there was as much genetic difference between these two groups as between either of them and all eukaryotes.
According to genetic data, although eukaryote groups such as plants, fungi, and animals may look different, they are more closely related to each other than they are to either the Eubacteria or Archaea. It was also found that the eukaryotes are more closely related to the Archaea than they are to the Eubacteria. Although the primacy of the Eubacteria-Archaea divide has been questioned, it has been upheld by subsequent research. There is no consensus on how many kingdoms exist in the classification scheme proposed by Woese.
Eukaryotic supergroups
In 2004, a review article by Simpson and Roger noted that the Protista were "a grab-bag for all eukaryotes that are not animals, plants or fungi". They held that only monophyletic groups should be accepted as formal ranks in a classification and that – while this approach had been impractical previously (necessitating "literally dozens of eukaryotic 'kingdoms) – it had now become possible to divide the eukaryotes into "just a few major groups that are probably all monophyletic".
On this basis, the diagram opposite (redrawn from their article) showed the real "kingdoms" (their quotation marks) of the eukaryotes. A classification which followed this approach was produced in 2005 for the International Society of Protistologists, by a committee which "worked in collaboration with specialists from many societies". It divided the eukaryotes into the same six "supergroups". The published classification deliberately did not use formal taxonomic ranks, including that of "kingdom".
In this system the multicellular animals (Metazoa) are descended from the same ancestor as both the unicellular choanoflagellates and the fungi which form the Opisthokonta. Plants are thought to be more distantly related to animals and fungi.
However, in the same year as the International Society of Protistologists' classification was published (2005), doubts were being expressed as to whether some of these supergroups were monophyletic, particularly the Chromalveolata, and a review in 2006 noted the lack of evidence for several of the six proposed supergroups.
, there is widespread agreement that the Rhizaria belong with the Stramenopiles and the Alveolata, in a clade dubbed the SAR supergroup, so that Rhizaria is not one of the main eukaryote groups.
Comparison of top level classification
Some authors have added non-cellular life to their classifications. This can create a "superdomain" called "Acytota", also called "Aphanobionta", of non-cellular life; with the other superdomain being "cytota" or cellular life. The eocyte hypothesis proposes that the eukaryotes emerged from a phylum within the archaea called the Thermoproteota (formerly known as eocytes or Crenarchaeota).
Viruses
The International Committee on Taxonomy of Viruses uses the taxonomic rank "kingdom" in the classification of viruses (with the suffix -virae); but this is beneath the top level classifications of realm and subrealm.
There is ongoing debate as to whether viruses can be included in the tree of life. The arguments against include the fact that they are obligate intracellular parasites that lack metabolism and are not capable of replication outside of a host cell. Another argument is that their placement in the tree would be problematic, since it is suspected that viruses have various evolutionary origins, and they have a penchant for harvesting nucleotide sequences from their hosts.
On the other hand, there are arguments in favor of their inclusion.
One of these comes from the discovery of unusually large and complex viruses, such as Mimivirus, that possess typical cellular genes.
| Biology and health sciences | Genetics and taxonomy | null |
53686 | https://en.wikipedia.org/wiki/Archimedes%27%20screw | Archimedes' screw | The Archimedes' screw, also known as the Archimedean screw, hydrodynamic screw, water screw or Egyptian screw, is one of the earliest hydraulic machines named after Greek mathematician Archimedes who first described it around 234 BC, although the device had been used in Ancient Egypt. It is a reversible hydraulic machine, and there are several examples of Archimedes screw installations where the screw can operate at different times as either pump or generator, depending on needs for power and watercourse flow.
As a machine used for lifting water from a low-lying body of water into irrigation ditches, water is lifted by turning a screw-shaped surface inside a pipe. In the modern world, Archimedes screw pumps are widely used in wastewater treatment plants and for dewatering low-lying regions. Run in reverse, Archimedes screw turbines act as a new form of small hydroelectric powerplant that can be applied even in low head sites. Such generators operate in a wide range of flows (0.01 to 14.5 ) and heads (0.1 m to 10 m), including low heads and moderate flow rates that is not ideal for traditional turbines and not occupied by high performance technologies.
History
Earliest records
The screw pump is the oldest positive displacement pump. The first records of a water screw, or screw pump, date back to Hellenistic Egypt before the 3rd century BC. The Egyptian screw, used to lift water from the Nile, was composed of tubes wound round a cylinder; as the entire unit rotates, water is lifted within the spiral tube to the higher elevation. A later screw pump design from Egypt had a spiral groove cut on the outside of a solid wooden cylinder and then the cylinder was covered by boards or sheets of metal closely covering the surfaces between the grooves.
Some researchers have proposed this device was used to irrigate the Hanging Gardens of Babylon, one of the Seven Wonders of the Ancient World. A cuneiform inscription of Assyrian King Sennacherib (704–681 BC) has been interpreted by Stephanie Dalley to describe casting water screws in bronze some 350 years earlier. This is consistent with Greek historian Strabo, who describes the Hanging Gardens as irrigated by screws.
Archimedes' role
The screw pump was later introduced from Hellenistic Egypt to Greece. It was described by Archimedes, on the occasion of his visit to Egypt, circa 234 BC. This tradition may reflect only that the apparatus was unknown to the Greeks before Hellenistic times. Athenaeus of Naucratis quotes a certain Moschion in a description on how Hiero II of Syracuse commissioned the design of the Syracusia, a luxury ship which would be a display of naval power. It is said to have been the largest ship built in classical antiquity and was launched by Archimedes who designed device with a revolving screw-shaped blade inside a cylinder to remove any potential water leaking through the hull. Archimedes' screw was turned by hand, and could also be used to transfer water from a low-lying body of water into irrigation canals.
Archimedes never claimed credit for its invention, but it was attributed to him 200 years later by Diodorus, who believed that Archimedes invented the screw pump in Egypt. Depictions of Greek and Roman water screws show them being powered by a human treading on the outer casing to turn the entire apparatus as one piece, which would require that the casing be rigidly attached to the screw.
Development and modern use
German engineer Konrad Kyeser equipped the Archimedes screw with a crank mechanism in his Bellifortis (1405). This mechanism quickly replaced the ancient practice of working the pipe by treading. The world's first seagoing steamship driven by a screw propeller was the SS Archimedes, which was launched in 1839 and named in honor of Archimedes and his work on the screw. Developments in maritime transport occurred over the next 180 years from the Fawcett, Preston and Company double blade design and patents by Sharrow Marine to address rotary propulsion and flow control on boating vessels through loop propellers. Electricity generation through hydropower pumps such as the Meriden project operated by New England Hydropower also uses Archimedes screw to direct water into the top, rather than the bottom, of the screw which forces it to rotate.
Archimedes screws are used in sewage treatment plants because they cope well with varying rates of flow and with suspended solids. Screw turbines (ASTs) are a new form of generator for small hydroelectric powerplants that could be applied even in low-head sites. The low rotation speed of ASTs reduces negative impacts on aquatic life and fish. This technology is used primarily at fish hatcheries to lift fish safely from ponds and transport them to another location. An Archimedes screw was used in the successful 2001 stabilization of the Leaning Tower of Pisa. Small amounts of subsoil saturated by groundwater were removed from far below the north side of the tower, and the weight of the tower itself corrected the lean.
Other inventions using Archimedes screws include the auger conveyor in a snow blower, grain elevator, concrete mixer and chocolate fountain.
Design
The Archimedes screw consists of a screw (a helical surface surrounding a central cylindrical shaft) inside a hollow pipe. The screw is usually turned by windmill, manual labor, cattle, or by modern means, such as a motor. As the shaft turns, the bottom end scoops up a volume of water. This water is then pushed up the tube by the rotating helicoid until it pours out from the top of the tube.
The contact surface between the screw and the pipe does not need to be perfectly watertight, as long as the amount of water being scooped with each turn is large compared to the amount of water leaking out of each section of the screw per turn. If water from one section leaks into the next lower one, it will be transferred upwards by the next segment of the screw.
In some designs, the screw is fused to the casing and they both rotate together, instead of the screw turning within a stationary casing. The screw could be sealed to the casing with pitch resin or other adhesive, or the screw and casing could be cast together as a single piece in bronze.
The design of the everyday Greek and Roman water screw, in contrast to the heavy bronze device of Sennacherib, with its problematic drive chains, has a powerful simplicity. A double or triple helix was built of wood strips (or occasionally bronze sheeting) around a heavy wooden pole. A cylinder was built around the helices using long, narrow boards fastened to their periphery and waterproofed with pitch.
Studies show that the volume of flow passes through Archimedes screws is a function of inlet depth, diameter and rotation speed of the screw. Therefore, the following analytical equation could be used to design Archimedes screws:
where is in and:
: Rotation speed of the Archimedes screw (rad/s)
: Volumetric flow rate
Based on the common standards that the Archimedes screw designers use this analytical equation could be simplified as:
The value of η could simply determinate using the graph or graph. By determination of , other design parameters of Archimedes screws can be calculated using a step-by-step analytical method.
Variants
A screw conveyor is a similar device which transports bulk materials such as powders and cereal grains. It is contained within a tube and turned by a motor to deliver material from one end of the conveyor to the other and particularly suitable for transport of granular materials such as plastic granules used in injection moulding. It may also be used to transport liquids. In industrial control applications, the conveyor may be used as a rotary feeder or variable rate feeder to deliver a measured rate or quantity of material into a process.
A variant of the Archimedes screw can also be found in some injection moulding machines, die casting machines and extrusion of plastics, which employ a screw of decreasing pitch to compress and melt the material. It is also used in a rotary-screw air compressor. On a much larger scale, Archimedes's screws of decreasing pitch are used for the compaction of waste material.
Reverse action
If water is fed into the top of an Archimedes screw, it will force the screw to rotate. The rotating shaft can then be used to drive an electric generator. Such an installation has the same benefits as using the screw for pumping: the ability to handle very dirty water and widely varying rates of flow at high efficiency. Settle Hydro and Torrs Hydro are two reverse screw micro hydro schemes operating in England. The screw works well as a generator at low heads, commonly found in English rivers, including the Thames, powering Windsor Castle.
| Technology | Agricultural tools | null |
53694 | https://en.wikipedia.org/wiki/Pig%20iron | Pig iron | Pig iron, also known as crude iron, is an intermediate good used by the iron industry in the production of steel. It is developed by smelting iron ore in a blast furnace. Pig iron has a high carbon content, typically 3.8–4.7%, along with silica and other dross, which makes it brittle and not useful directly as a material except for limited applications.
Etymology
The traditional shape of the molds used for pig iron ingots is a branching structure formed in sand, with many individual ingots at right angles to a central channel or "runner", resembling a litter of piglets being nursed by a sow. When the metal had cooled and hardened, the smaller ingots (the "pigs") were simply broken from the runner (the "sow"), hence the name "pig iron". As pig iron is intended for remelting, the uneven size of the ingots and the inclusion of small amounts of sand are insignificant issues when compared to the ease of casting and handling.
History
The Chinese were already making pig iron during the later Zhou dynasty (which ended in 256 BC). Furnaces such as Lapphyttan in Sweden may date back as far back as the 12th century; and some in the County of Mark dating back to the 13th century, which is now part of Westphalia, Germany. It remains to be established whether these northern European developments were derived from the Chinese ones. Wagner has postulated a possible link via Persian contacts with China along the Silk Road and Viking contacts with Persia, but there is a chronological gap between the Viking period and Lapphyttan.
Smelting and producing wrought iron were known in ancient Europe and the Middle East, but it was produced in bloomeries by direct reduction. Small prills of pig iron dispersed in slag are produced in all iron furnaces, but the operator of a bloomery had to avoid conditions causing the phase transition of the iron into liquid in the furnace, as the prill globules or any resulting pig iron are not malleable so can't be hammered in a single piece. Alternatively, decarburizing the pig iron into steel was an extremely tedious process using medieval technology, so in Europe before the Middle Ages the prills were discarded with the slag.
Uses
Traditionally, pig iron was worked into wrought iron in finery forges, later puddling furnaces, and more recently, into steel. In these processes, pig iron is melted and a strong current of air is directed over it while it is stirred or agitated. This causes the dissolved impurities (such as silicon) to be thoroughly oxidized. An intermediate product of puddling is known as refined pig iron, finers metal, or refined iron.
Pig iron can also be used to produce gray iron. This is achieved by remelting pig iron, often along with substantial quantities of steel and scrap iron, removing undesirable contaminants, adding alloys, and adjusting the carbon content. Ductile iron can also be produced using certain high purity grades of pig iron; depending on the grade of ductile iron being produced, the pig irons chosen may be low in the elements silicon, manganese, sulfur and phosphorus. High purity pig iron is used to dilute any elements in a ductile iron charge which may be harmful to the ductile iron process (except carbon).
Modern uses
Pig iron was historically poured directly out of the bottom of the blast furnace through a trough into a ladle car for transfer to the steel mill in mostly liquid form; in this state, the pig iron was referred to as hot metal. The hot metal was then poured into a steelmaking vessel to produce steel, typically an electric arc furnace, induction furnace or basic oxygen furnace, where the excess carbon is burned off and the alloy composition controlled. Earlier processes for this included the finery forge, the puddling furnace, the Bessemer process, and the open hearth furnace.
Modern steel mills and direct-reduction iron plants transfer the molten iron to a ladle for immediate use in the steel making furnaces or cast it into pigs on a pig-casting machine for reuse or resale. Modern pig casting machines produce stick pigs, which break into smaller piglets at discharge.
| Physical sciences | Iron alloys | Chemistry |
53696 | https://en.wikipedia.org/wiki/Division%20%28mathematics%29 | Division (mathematics) | Division is one of the four basic operations of arithmetic. The other operations are addition, subtraction, and multiplication. What is being divided is called the dividend, which is divided by the divisor, and the result is called the quotient.
At an elementary level the division of two natural numbers is, among other possible interpretations, the process of calculating the number of times one number is contained within another. For example, if 20 apples are divided evenly between 4 people, everyone receives 5 apples (see picture). However, this number of times or the number contained (divisor) need not be integers.
The division with remainder or Euclidean division of two natural numbers provides an integer quotient, which is the number of times the second number is completely contained in the first number, and a remainder, which is the part of the first number that remains, when in the course of computing the quotient, no further full chunk of the size of the second number can be allocated. For example, if 21 apples are divided between 4 people, everyone receives 5 apples again, and 1 apple remains.
For division to always yield one number rather than an integer quotient plus a remainder, the natural numbers must be extended to rational numbers or real numbers. In these enlarged number systems, division is the inverse operation to multiplication, that is means , as long as is not zero. If , then this is a division by zero, which is not defined. In the 21-apples example, everyone would receive 5 apple and a quarter of an apple, thus avoiding any leftover.
Both forms of division appear in various algebraic structures, different ways of defining mathematical structure. Those in which a Euclidean division (with remainder) is defined are called Euclidean domains and include polynomial rings in one indeterminate (which define multiplication and addition over single-variabled formulas). Those in which a division (with a single result) by all nonzero elements is defined are called fields and division rings. In a ring the elements by which division is always possible are called the units (for example, 1 and −1 in the ring of integers). Another generalization of division to algebraic structures is the quotient group, in which the result of "division" is a group rather than a number.
Introduction
The simplest way of viewing division is in terms of quotition and partition: from the quotition perspective, means the number of 5s that must be added to get 20. In terms of partition, means the size of each of 5 parts into which a set of size 20 is divided. For example, 20 apples divide into five groups of four apples, meaning that "twenty divided by five is equal to four". This is denoted as , or . In the example, 20 is the dividend, 5 is the divisor, and 4 is the quotient.
Unlike the other basic operations, when dividing natural numbers there is sometimes a remainder that will not go evenly into the dividend; for example, leaves a remainder of 1, as 10 is not a multiple of 3. Sometimes this remainder is added to the quotient as a fractional part, so is equal to or , but in the context of integer division, where numbers have no fractional part, the remainder is kept separately (or exceptionally, discarded or rounded). When the remainder is kept as a fraction, it leads to a rational number. The set of all rational numbers is created by extending the integers with all possible results of divisions of integers.
Unlike multiplication and addition, division is not commutative, meaning that is not always equal to . Division is also not, in general, associative, meaning that when dividing multiple times, the order of division can change the result. For example, , but (where the use of parentheses indicates that the operations inside parentheses are performed before the operations outside parentheses).
Division is traditionally considered as left-associative. That is, if there are multiple divisions in a row, the order of calculation goes from left to right:
Division is right-distributive over addition and subtraction, in the sense that
This is the same for multiplication, as . However, division is not left-distributive, as
For example but
This is unlike the case in multiplication, which is both left-distributive and right-distributive, and thus distributive.
Notation
Division is often shown in algebra and science by placing the dividend over the divisor with a horizontal line, also called a fraction bar, between them. For example, "a divided by b" can be written as:
which can also be read out loud as "divide a by b" or "a over b". A way to express division all on one line is to write the dividend (or numerator), then a slash, then the divisor (or denominator), as follows:
This is the usual way of specifying division in most computer programming languages, since it can easily be typed as a simple sequence of ASCII characters. (It is also the only notation used for quotient objects in abstract algebra.) Some mathematical software, such as MATLAB and GNU Octave, allows the operands to be written in the reverse order by using the backslash as the division operator:
A typographical variation halfway between these two forms uses a solidus (fraction slash), but elevates the dividend and lowers the divisor:
Any of these forms can be used to display a fraction. A fraction is a division expression where both dividend and divisor are integers (typically called the numerator and denominator), and there is no implication that the division must be evaluated further. A second way to show division is to use the division sign (÷, also known as obelus though the term has additional meanings), common in arithmetic, in this manner:
This form is infrequent except in elementary arithmetic. ISO 80000-2-9.6 states it should not be used. This division sign is also used alone to represent the division operation itself, as for instance as a label on a key of a calculator. The obelus was introduced by Swiss mathematician Johann Rahn in 1659 in Teutsche Algebra. The ÷ symbol is used to indicate subtraction in some European countries, so its use may be misunderstood.
In some non-English-speaking countries, a colon is used to denote division:
This notation was introduced by Gottfried Wilhelm Leibniz in his 1684 Acta eruditorum. Leibniz disliked having separate symbols for ratio and division. However, in English usage the colon is restricted to expressing the related concept of ratios.
Since the 19th century, US textbooks have used or to denote a divided by b, especially when discussing long division. The history of this notation is not entirely clear because it evolved over time.
Computing
Manual methods
Division is often introduced through the notion of "sharing out" a set of objects, for example a pile of lollies, into a number of equal portions. Distributing the objects several at a time in each round of sharing to each portion leads to the idea of 'chunking' a form of division where one repeatedly subtracts multiples of the divisor from the dividend itself.
By allowing one to subtract more multiples than what the partial remainder allows at a given stage, more flexible methods, such as the bidirectional variant of chunking, can be developed as well.
More systematically and more efficiently, two integers can be divided with pencil and paper with the method of short division, if the divisor is small, or long division, if the divisor is larger. If the dividend has a fractional part (expressed as a decimal fraction), one can continue the procedure past the ones place as far as desired. If the divisor has a fractional part, one can restate the problem by moving the decimal to the right in both numbers until the divisor has no fraction, which can make the problem easier to solve (e.g., 10/2.5 = 100/25 = 4).
Division can be calculated with an abacus.
Logarithm tables can be used to divide two numbers, by subtracting the two numbers' logarithms, then looking up the antilogarithm of the result.
Division can be calculated with a slide rule by aligning the divisor on the C scale with the dividend on the D scale. The quotient can be found on the D scale where it is aligned with the left index on the C scale. The user is responsible, however, for mentally keeping track of the decimal point.
By computer
Modern calculators and computers compute division either by methods similar to long division, or by faster methods; see Division algorithm.
In modular arithmetic (modulo a prime number) and for real numbers, nonzero numbers have a multiplicative inverse. In these cases, a division by may be computed as the product by the multiplicative inverse of . This approach is often associated with the faster methods in computer arithmetic.
Division in different contexts
Euclidean division
Euclidean division is the mathematical formulation of the outcome of the usual process of division of integers. It asserts that, given two integers, a, the dividend, and b, the divisor, such that b ≠ 0, there are unique integers q, the quotient, and r, the remainder, such that a = bq + r and 0 ≤ r < , where denotes the absolute value of b.
Of integers
Integers are not closed under division. Apart from division by zero being undefined, the quotient is not an integer unless the dividend is an integer multiple of the divisor. For example, 26 cannot be divided by 11 to give an integer. Such a case uses one of five approaches:
Say that 26 cannot be divided by 11; division becomes a partial function.
Give an approximate answer as a floating-point number. This is the approach usually taken in numerical computation.
Give the answer as a fraction representing a rational number, so the result of the division of 26 by 11 is (or as a mixed number, so ) Usually the resulting fraction should be simplified: the result of the division of 52 by 22 is also . This simplification may be done by factoring out the greatest common divisor.
Give the answer as an integer quotient and a remainder, so To make the distinction with the previous case, this division, with two integers as result, is sometimes called Euclidean division, because it is the basis of the Euclidean algorithm.
Give the integer quotient as the answer, so This is the floor function applied to case 2 or 3. It is sometimes called integer division, and denoted by "//".
Dividing integers in a computer program requires special care. Some programming languages treat integer division as in case 5 above, so the answer is an integer. Other languages, such as MATLAB and every computer algebra system return a rational number as the answer, as in case 3 above. These languages also provide functions to get the results of the other cases, either directly or from the result of case 3.
Names and symbols used for integer division include div, /, \, and %. Definitions vary regarding integer division when the dividend or the divisor is negative: rounding may be toward zero (so called T-division) or toward −∞ (F-division); rarer styles can occur – see modulo operation for the details.
Divisibility rules can sometimes be used to quickly determine whether one integer divides exactly into another.
Of rational numbers
The result of dividing two rational numbers is another rational number when the divisor is not 0. The division of two rational numbers p/q and r/s can be computed as
All four quantities are integers, and only p may be 0. This definition ensures that division is the inverse operation of multiplication.
Of real numbers
Division of two real numbers results in another real number (when the divisor is nonzero). It is defined such that a/b = c if and only if a = cb and b ≠ 0.
Of complex numbers
Dividing two complex numbers (when the divisor is nonzero) results in another complex number, which is found using the conjugate of the denominator:
This process of multiplying and dividing by is called 'realisation' or (by analogy) rationalisation. All four quantities p, q, r, s are real numbers, and r and s may not both be 0.
Division for complex numbers expressed in polar form is simpler than the definition above:
Again all four quantities p, q, r, s are real numbers, and r may not be 0.
Of polynomials
One can define the division operation for polynomials in one variable over a field. Then, as in the case of integers, one has a remainder. See Euclidean division of polynomials, and, for hand-written computation, polynomial long division or synthetic division.
Of matrices
One can define a division operation for matrices. The usual way to do this is to define , where denotes the inverse of B, but it is far more common to write out explicitly to avoid confusion. An elementwise division can also be defined in terms of the Hadamard product.
Left and right division
Because matrix multiplication is not commutative, one can also define a left division or so-called backslash-division as . For this to be well defined, need not exist, however does need to exist. To avoid confusion, division as defined by is sometimes called right division or slash-division in this context.
With left and right division defined this way, is in general not the same as , nor is the same as . However, it holds that and .
Pseudoinverse
To avoid problems when and/or do not exist, division can also be defined as multiplication by the pseudoinverse. That is, and , where and denote the pseudoinverses of and .
Abstract algebra
In abstract algebra, given a magma with binary operation ∗ (which could nominally be termed multiplication), left division of b by a (written ) is typically defined as the solution x to the equation , if this exists and is unique. Similarly, right division of b by a (written ) is the solution y to the equation . Division in this sense does not require ∗ to have any particular properties (such as commutativity, associativity, or an identity element). A magma for which both and exist and are unique for all a and all b (the Latin square property) is a quasigroup. In a quasigroup, division in this sense is always possible, even without an identity element and hence without inverses.
"Division" in the sense of "cancellation" can be done in any magma by an element with the cancellation property. Examples include matrix algebras, quaternion algebras, and quasigroups. In an integral domain, where not every element need have an inverse, division by a cancellative element a can still be performed on elements of the form ab or ca by left or right cancellation, respectively. If a ring is finite and every nonzero element is cancellative, then by an application of the pigeonhole principle, every nonzero element of the ring is invertible, and division by any nonzero element is possible. To learn about when algebras (in the technical sense) have a division operation, refer to the page on division algebras. In particular Bott periodicity can be used to show that any real normed division algebra must be isomorphic to either the real numbers R, the complex numbers C, the quaternions H, or the octonions O.
Calculus
The derivative of the quotient of two functions is given by the quotient rule:
Division by zero
Division of any number by zero in most mathematical systems is undefined, because zero multiplied by any finite number always results in a product of zero. Entry of such an expression into most calculators produces an error message. However, in certain higher level mathematics division by zero is possible by the zero ring and algebras such as wheels. In these algebras, the meaning of division is different from traditional definitions.
Calculator
Enter two numbers to find their quotient:
/ =
| Mathematics | Basics | null |
53700 | https://en.wikipedia.org/wiki/Universal%20Product%20Code | Universal Product Code | The Universal Product Code (UPC or UPC code) is a barcode symbology that is used worldwide for tracking trade items in stores.
The chosen symbology has bars (or spaces) of exactly 1, 2, 3, or 4 units wide each; each decimal digit to be encoded consists of two bars and two spaces chosen to have a total width of 7 units, in both an "even" and an "odd" parity form, which enables being scanned in either direction. Special "guard patterns" (3 or 5 units wide, not encoding a digit) are intermixed to help decoding.
A UPC (technically, a UPC-A) consists of 12 digits that are uniquely assigned to each trade item. The international GS1 organisation assigns the digits used for both the UPC and the related International Article Number (EAN) barcode. UPC data structures are a component of Global Trade Item Numbers (GTINs) and follow the global GS1 specification, which is based on international standards. Some retailers, such as clothing and furniture, do not use the GS1 system, instead using other barcode symbologies or article number systems. Some retailers use the EAN/UPC barcode symbology, but do not use a GTIN for products sold only in their own stores.
Research indicates that the adoption and diffusion of the UPC stimulated innovation and contributed to the growth of international retail supply chains.
History
Wallace Flint proposed an automated checkout system in 1932 using punched cards. Bernard Silver and Norman Joseph Woodland, a graduate student from Drexel Institute of Technology, developed a bull's-eye-style code and applied for the patent in 1949.
In the 1960s and early 1970s, railroads in North America experimented with multicolor bar codes for tracking railcars, but this system was eventually abandoned and replaced with a radio-based system called Automatic Equipment Identification (AEI).
In 1973, a group of trade associations from the grocery industry formed the Uniform Product Code Council (UPCC) which, with the help of consultants Larry Russell and Tom Wilson of McKinsey & Company, defined the numerical format that formed the basis of the Uniform Product Code. Technology firms including Charegon, IBM, Litton-Zellweger, Pitney Bowes-Alpex, Plessey-Anker, RCA, Scanner Inc., Singer, and Dymo Industries/Data General, put forward alternative proposals for symbol representations to the council. The Symbol Selection Committee finally chose to implement the IBM proposal designed by George J. Laurer, but with a slight modification to the font in the human readable area.
The first UPC-marked item ever to be scanned at a retail checkout was a 10-pack (50 sticks) of Wrigley's Juicy Fruit chewing gum, purchased at the Marsh supermarket in Troy, Ohio, at 8:01 a.m. on 26 June 1974. The NCR cash register rang up 67 cents.
The shopping cart also contained other barcoded items but the gum was the first one picked up at the checkout. A facsimile of the gum packet went on display at the Smithsonian Institution's American history museum in Washington, D.C.
Murray Eden was a consultant on the team that created the Universal Product Code barcode. As Chairman of a committee of scientists at the Massachusetts Institute of Technology, he helped "select a symbol that would endure the inevitable rush of technology that lay ahead." He chose the font, and he came up with the idea to add numbers to the bottom, which is a fail-safe system, in case the barcode reader is not working correctly.
IBM proposal
Around late 1969, IBM at Research Triangle Park (RTP) in North Carolina assigned George Laurer to determine how to make a supermarket scanner and label. In late 1970, Heard Baumeister provided equations to calculate characters-per-inch achievable by two IBM bar codes, Delta A and Delta B. In February 1971, Baumeister joined Laurer.
Delta B compared bar widths to space width to code bits. This was extremely sensitive to ink spread, where improper levels of ink or pressure would cause both edges of a bar to spread outward or shrink in.
In mid 1971, William "Bill" Crouse invented a new bar code called Delta C.
It achieved four times the characters per inch as Delta B.
Delta C achieved its higher performance by only using leading to leading or trailing to trailing edges which was unaffected by uniform ink spread. The code provided best performance when it had a defined character set with a fixed reference distance that spanned most or preferably all the character.
In August 1971, Crouse joined the scanner effort. After several months they had made no progress. They were aware of the RCA bull's eye label that could be scanned with a simple straight line laser scanner, but a readable label was far too large.
Although Litton Industries proposed a bull's eye symbol cut in half to reduce the area, it was still too large and presented the same ink smear printing problems as the RCA symbol. The redundancy and checking ability were removed completely. They were also aware of the many proposals from around the world, none of which were feasible.
In the spring of 1972, Baumeister announced a breakthrough. He proposed a label with bars that were slightly longer than the distance across all bars that needed to be read in a single pass. This label could be scanned with a simple "X" scanner only slightly more complex than the straight line laser scanner. The next day Baumeister suggested if the label were split into two halves the bar lengths could be cut nearly in half.
These two proposals reduced the area from the bull's eye by one third and then one sixth. The image to the right shows the label proposed by Baumeister. He did not specify any specific bar code as that was well understood. Except for the bar coding and ten digits the UPC label today is his proposal. Shortly after that Baumeister transferred to another area of RTP.
Laurer proceeded to define the details of the label and write a proposal. N.J. Woodland was assigned as planner for the project and aided Laurer with writing his proposal.
Laurer's first attempt with a bar code used Delta B. The resulting label size was about six inches by three inches which was too large. Crouse suggested that Laurer use his Delta C bar code and provided a copy of his patent that had a sample alphanumeric character set and rules to generate other size alphabets. This reduced the label size to about .
Later Laurer asked Crouse for assistance in how the scanner could detect a label. Together they defined guard bars and a definition of how to detect the label. The guard bars also provided identification for half label discrimination and training bars for the scanner threshold circuits. Laurer had a complete label definition and proceeded to write his proposal.
Previously Crouse had an idea for a simple wand worn like a ring and bracelet. He decided to develop that wand to provide a demonstration of the label.
On 1 December 1972, IBM presented Laurer's proposal to the Super Market Committee in Rochester, Minnesota, the location where IBM would develop the scanner. During the presentation, Crouse gave a lab demonstration where he read UPC-like labels with his ring wand. In addition to reading regular labels, he read the large two-page centerfold label in the proposal booklet. He then turned to a page showing a photo of labeled items sitting on a table. The labels were small and flawed due to the resolution of the printed photo but the wand read many of them. This demonstration showed the robustness of the pure Delta C code. The proposal was accepted.
One month later, 1 January 1973 Crouse transferred back to IBM's Advanced Technology group, and Laurer remained with the full responsibility for the label.
Dymo Industries, makers of handheld printing devices insisted that the code be character independent, so that handheld printing devices could produce the bar code in store if the items were not bar-coded by the manufacturers. Dymo's proposal was accepted by IBM and incorporated in IBM's latest proposal.
It was decided that the two halves of the label should have a different set of numeric characters. The character set Laurer derived from the Delta C patent used seven printable increments or units where two bars and two spaces would be printed. This yielded twenty combinations of characters, but there were two pairs that when read by Delta C rules yielded the same code for the pair.
Since eighteen characters were not enough Laurer tried adding one unit to the character set. This yielded twenty-six Delta C characters which could provide the two sets of decimal characters but it also added fourteen percent to the width of the label and thereby the height. This would be a thirty percent increase in area or a label of . Laurer felt this was not acceptable.
Laurer returned to the original character set with twenty characters but four of those were two pairs with the same Delta C reading. He decided to use them all. To distinguish between the pairs he would measure one bar width in each of the pairs to distinguish them from each other. For each pair those bars would be one or two units wide.
Laurer did not apply Baumeister's equations to this set. He felt just one bar width measurement would not be too serious. As it turned out it would have required over fifty percent increase in width and height for an area increase of more than double. Laurer later admitted these four characters in each set were responsible for most of the scanner read errors.
David Savir, a mathematician, was given the task of proving the symbol could be printed and would meet the reliability requirements, and was most likely unaware of Baumeister's equations. He and Laurer added two more digits to the ten for error detection and correction.
Then they decided to add odd/even parity to the number of units filled with bars in each side. Odd/even parity is a technique used to detect any odd number of bit errors in a bit stream. They decided to use odd on one half and even on the other. This would provide additional indication of which half ticket was being read. This meant that every bar width had to be read accurately to provide a good reading. It also meant every space would also be known.
Requiring every bit width to be read precisely basically nullified the Delta C advantage except for the Delta C reference measurement. Only the strange character set and the size of the label remains as a shadow of the Delta C code. The size was still that calculated for pure Delta C. If the label size had been properly recalculated, taking into account the required bar width measurements the label would have been far too large to be acceptable.
Mechanical engineering and electronic circuit design commonly require worst case designs using known tolerances. Many engineers working with bar codes had little experience with such things and used somewhat intuitive methods. This was the cause of the poor performance of the Delta B code and quite likely the failure of RCA's bull's eye scanner.
The following table shows the workable labels, available in the early 1970s, with their sizes.
This is assuming a bull's eye with the same information and reliable readability.
Composition
Each UPC-A barcode consists of a scannable strip of black bars and white spaces above a sequence of 12 numerical digits. No alphabetic characters, symbols, or other characters of any kind can appear on a UPC-A barcode. There is a one-to-one correspondence between 12-digit number and strip of black bars and white spaces, i.e. there is only one way to represent each 12-digit number visually and there is only one way to represent each strip of black bars and white spaces numerically.
The scannable area of every UPC-A barcode follows the pattern SLLLLLLMRRRRRRE, where S (start), M (middle), and E (end) guard patterns are represented the same way on every UPC-A barcode and the L (left) and R (right) sections collectively represent the 12 numerical digits that make each UPC-A unique. The first digit L indicates a particular number system to be used by the following digits. The last digit R is an error detecting check digit, that allows some errors to be detected in scanning or manual entry. The guard patterns separate the two groups of six numerical digits and establish the timing.
UPC-A 042100005264 is equivalent to UPC-E 425261 with the "EOEEOO" parity pattern, which is defined by UPC-A number system 0 and UPC-A check digit 4.
Formatting
UPC-A barcodes can be printed at various densities to accommodate a variety of printing and scanning processes. The significant dimensional parameter is called x-dimension (width of single module element). The width of each bar (space) is determined by multiplying the x-dimension and the module width (1, 2, 3, or 4 units) of each bar (space). Since the guard patterns each include two bars, and each of the 12 digits of the UPC-A barcode consists of two bars and two spaces, all UPC-A barcodes consist of exactly (3 × 2) + (12 × 2) = 30 bars, of which 6 represent guard patterns and 24 represent numerical digits.
The x-dimension for the UPC-A at the nominal size is 0.33 mm (0.013"). Nominal symbol height for UPC-A is 25.9 mm (1.02"). The bars forming the S (start), M (middle), and E (end) guard patterns, are extended downwards by 5 times x-dimension, with a resulting nominal symbol height of 27.55 mm (1.08"). This also applies to the bars of the first and last numerical digit of UPC-A barcode. UPC-A can be reduced or magnified anywhere from 80% to 200%.
A quiet zone, with a width of at least 9 times the x-dimension, must be present on each side of the scannable area of the UPC-A barcode. For a GTIN-12 number encoded in a UPC-A barcode, the first and last digits of the human-readable interpretation are always placed outside the symbol in order to indicate the quiet zones that are necessary for UPC barcode scanners to work properly.
Encoding
The UPC-A barcode is visually represented by strips of bars and spaces that encode the UPC-A 12-digit number. Each digit is represented by a unique pattern of 2 bars and 2 spaces. The bars and spaces are variable width, i.e. 1, 2, 3, or 4 modules wide. The total width for a digit is always 7 modules; consequently, UPC-A 12-digit number requires a total of .
A complete UPC-A is 95 modules wide: 84 modules for the digits (L and R sections) combined with 11 modules for the S (start), M (middle), and E (end) guard patterns. The S (start) and E (end) guard patterns are 3 modules wide and use the pattern bar-space-bar, where each bar and space is one module wide. The M (middle) guard pattern is 5 modules wide and uses the pattern space-bar-space-bar-space, where each bar and space is also one module wide. In addition, a UPC-A symbol requires a quiet zone (extra space of 9 modules wide) before the S (start) and after the E (end) guard patterns.
The UPC-A's left-hand side digits (the digits to the left of the M (middle) guard pattern) have odd parity, which means the total width of the black bars is an odd number of modules. On the contrary, the right-hand side digits have even parity. Consequently, a UPC scanner can determine whether it is scanning a symbol from left-to-right or from right-to-left (the symbol is upside-down). After seeing a S (start) or E (end) guard pattern (they are the same, bar-space-bar, whichever direction they are read), the scanner will first see odd parity digits, if scanning left-to-right, or even parity digits, if scanning right-to-left. With the parity/direction information, an upside-down symbol will not confuse the scanner. When confronted with an upside-down symbol, the scanner may simply ignore it (many scanners alternate left-to-right and right-to-left scans, so they will read the symbol on a subsequent pass) or recognize the digits and put them in the right order. There is another property in the digit encoding. The right-hand side digits are the optical inverse of the left-hand side digits, i.e. black bars are turned into white spaces and vice versa. For example, the left-hand side "4" is , meanwhile the right-hand side "4" is .
Numbering
The number of UPC-A and UPC-E barcodes are limited by the standards used to create them.
UPC-A
UPC-E
Number system digit
Below is a description of all possible number systems with the corresponding 12-digit UPC-A numbering scheme LLLLLLRRRRRR, where L denotes the numbering system digit and R the check digit.
0–1, 6–9 For most products. The LLLLL digits are the manufacturer code (assigned by local GS1 organization), and the RRRRR digits are the product code.
2 Reserved for local use (store/warehouse), for items sold by variable weight. Variable-weight items, such as meats, fresh fruits, or vegetables, are assigned an item number by the store, if they are packaged there. In this case, the LLLLL is the item number, and the RRRRR is either the weight or the price, with the first R determining which (0 for weight).
3 Drugs by National Drug Code (NDC) number. Pharmaceuticals in the U.S. use the middle 10 digits of the UPC as their NDC number. Though usually only over-the-counter drugs are scanned at point of sale, NDC-based UPCs are used on prescription drug packages and surgical products and, in this case, are commonly called UPN Codes.
4 Reserved for local use (store/warehouse), often for loyalty cards or store coupons.
5 Coupons. The LLLLL digits are digits 2-6 of the product's UPC prefix, the next three RRR are a family code (set by manufacturer or supplied by the coupon clearing house), and the next two RR are a value code (according to the GS1 value code table), which determines the amount of the discount. These coupons can be doubled or tripled.
Check digit calculation
The UPC includes a check digit to detect common data entry errors. For example, UPC-A codes choose the check digit to satisfy the check digit equation:
If an entered code does not satisfy the equation, then it is not a valid UPC-A.
The UPC-A check digit may be calculated as follows:
Sum the digits at odd-numbered positions (first, third, fifth,..., eleventh).
Multiply the result by 3.
Add the digit sum at even-numbered positions (second, fourth, sixth,..., tenth) to the result.
Find the result modulo 10 (i.e. the remainder, when divided by 10) and call it .
If is zero, then the check digit is 0; otherwise the check digit is .
For example, in a UPC-A barcode "03600029145x12", where is the unknown check digit, may be calculated by:
Sum the odd-numbered digits (0 + 6 + 0 + 2 + 1 + 5 = 14).
Multiply the result by 3 (14 × 3 = 42).
Add the even-numbered digits (42 + (3 + 0 + 0 + 9 + 4) = 58).
Find the result modulo 10 (58 mod 10 = 8 = M).
If is not 0, subtract from 10 ().
Thus, the check digit is 2.
The check digit equation is selected to have reasonable error detection properties (see Luhn algorithm).
UPC-A can detect 100% of single digit errors.
A single digit error means exactly one digit is wrong. Let the difference modulo 10 of the erroneous digit and the correct digit be . The value of cannot be zero because that means the digits are the same, but can be any other value in {1, 2, 3, 4, 5, 6, 7, 8, 9}. If the error digit is in an odd position (weight 1), the left hand side of check digit equation changes by and the equivalence is no longer zero. If the error digit is in an even position (weight 3), then the left hand side changes by , but that change is also nonzero modulo 10, so the check digit equation is not satisfied.
UPC-A can detect about 89% of transposition errors. Specifically, if and only if the difference between two adjacent digits is 5, the UPC-A can't detect their transposition.
If 2 neighboring digits are transposed, then one of the digits will be weighted by 1, and the other digit will be weighted by 3, where is the difference between the two digits. If the digits were in their correct order, they would contribute
to the left hand side of the check digit equation. In the transposed order, they contribute
.
to the LHS. Subtracting the two contributions gives how much they change the LHS:
An error will be detected as long as the modular change is nonzero; if modulo 10, then the change will not be detected. Consequently, only when the character difference will an error be undetected (when the degenerate "transposition" is not an error).
Next consider how often a transposition has a distance of 5.
Here is the Table of d-transpositions for UPC-A barcodes, where
{|class="wikitable"
|+ Table of d-transpositions for UPC-A barcodes
!
! 0 !! 1 !! 2 !! 3 !! 4 !! 5 !! 6 !! 7 !! 8 !! 9
|-
! 1
| 0 0 || 0 1 || 0 2 || 0 3 || 0 4 || 0 5 || 0 6 || 0 7 || 0 8 || 0 9
|-
! 2
| 1 1 || 1 2 || 1 3 || 1 4 || 1 5 || 1 6 || 1 7 || 1 8 || 1 9
|-
! 3
| 2 2 || 2 3 || 2 4 || 2 5 || 2 6 || 2 7 || 2 8 || 2 9
|-
! 4
| 3 3 || 3 4 || 3 5 || 3 6 || 3 7 || 3 8 || 3 9
|-
! 5
| 4 4 || 4 5 || 4 6 || 4 7 || 4 8 || 4 9
|-
! 6
| 5 5 || 5 6 || 5 7 || 5 8 || 5 9
|-
! 7
| 6 6 || 6 7 || 6 8 || 6 9
|-
! 8
| 7 7 || 7 8 || 7 9
|-
! 9
| 8 8 || 8 9
|-
! 10
| 9 9
|-
! Sum
| 10 || 18 || 16 || 14 || 12 || 10 || 8 || 6 || 4 || 2
|}
Row Sum contains the number of d-transpositions, therefore the proportion of non-detectable transposition errors is (ignoring the transpositions where ):
▯
Variations
UPC in its most common usage technically refers to UPC-A.
Other variants of the UPC exist:
UPC-B is a 12-digit version of UPC with no check digit, developed for the National Drug Code (NDC) and National Health Related Items Code. It has 11 digits plus a 1-digit product code, and is not in common use.
UPC-C is a 12-digit code with a product code and a check digit; not in common use.
UPC-D is a variable length code (12 digits or more) with the 12th digit being the check digit. These versions are not in common use.
UPC-E is a 6-digit code, that has its equivalent in UPC-A 12-digit code with number system 0 or 1.
UPC-2 is a 2-digit supplement to the UPC used to indicate the edition of a magazine or periodical.
UPC-5 is a 5-digit supplement to the UPC used to indicate suggested retail price for books.
UPC-E
To allow the use of UPC barcodes on smaller packages, where a full 12-digit barcode may not fit, a zero-suppressed version of UPC was developed, called UPC-E, in which the number system digit, all trailing zeros in the manufacturer code, and all leading zeros in the product code, are suppressed (omitted). This symbology differs from UPC-A in that it only uses a 6-digit code, does not use M (middle) guard pattern, and the E (end) guard pattern is formed as space-bar-space-bar-space-bar, i.e. UPC-E barcode follows the pattern SDDDDDDE. The way in which a 6-digit UPC-E relates to a 12-digit UPC-A, is determined by UPC-E numerical pattern and UPC-E parity pattern. It can only correspond to UPC-A number system 0 or 1, the value of which, along with the UPC-A check digit, determines the UPC-E parity pattern of the encoding. With the manufacturer code digits represented by M's, and product code digits by P's, then:
For example, a UPC-E 654321 may correspond to the UPC-A 065100004327 or 165100004324, depending on the UPC-E parity pattern of the encoded digits, as described next:
UPC-E 654321 with "EOEOEO" parity pattern (UPC-A 065100004327) would be encoded as
1-1-1 4-1-1-1 1-2-3-1 2-3-1-1 1-4-1-1 2-2-1-2 2-2-2-1 1-1-1-1-1-1.
The barcode would look like this:
EAN-13
The EAN-13 was developed as a superset of UPC-A, adding an extra digit to the beginning of every UPC-A number. This expanded the number of unique values theoretically possible by ten times to 1 trillion. EAN-13 barcodes also indicate the country in which the company that sells the product is based (which may or may not be the same as the country in which the good is manufactured). The three leading digits of the code determine this, according to the GS1 country codes. Every UPC-A code can be easily converted to the equivalent EAN-13 code by prepending 0 digit to the UPC-A code. This does not change the check digit. All point-of-sale systems can now understand both equally.
EAN-8 is an 8-digit variation of the EAN barcode.
UPC usage notes:
All products marked with an EAN will be accepted in North America currently, in addition to those already marked with a UPC.
Products with an existing UPC do not have to be re-marked with an EAN.
In North America, the EAN adds 30% more codes, mainly by adding digits 10 through 12 to the UPC digits 00 through 09. This is a powerful incentive to phase out the UPC.
| Technology | Software development: General | null |
53702 | https://en.wikipedia.org/wiki/Boltzmann%20constant | Boltzmann constant | The Boltzmann constant ( or ) is the proportionality factor that relates the average relative thermal energy of particles in a gas with the thermodynamic temperature of the gas. It occurs in the definitions of the kelvin (K) and the gas constant, in Planck's law of black-body radiation and Boltzmann's entropy formula, and is used in calculating thermal noise in resistors. The Boltzmann constant has dimensions of energy divided by temperature, the same as entropy and heat capacity. It is named after the Austrian scientist Ludwig Boltzmann.
As part of the 2019 revision of the SI, the Boltzmann constant is one of the seven "defining constants" that have been defined so as to have exact finite decimal values in SI units. They are used in various combinations to define the seven SI base units. The Boltzmann constant is defined to be exactly joules per kelvin. Correspondingly, the SI units for temperature and energy are calibrated to one another so that kelvin = joules.
Roles of the Boltzmann constant
Macroscopically, the ideal gas law states that, for an ideal gas, the product of pressure and volume is proportional to the product of amount of substance and absolute temperature :
where is the molar gas constant (). Introducing the Boltzmann constant as the gas constant per molecule (NA being the Avogadro constant) transforms the ideal gas law into an alternative form:
where is the number of molecules of gas.
Role in the equipartition of energy
Given a thermodynamic system at an absolute temperature , the average thermal energy carried by each microscopic degree of freedom in the system is (i.e., about , or , at room temperature). This is generally true only for classical systems with a large number of particles, and in which quantum effects are negligible.
In classical statistical mechanics, this average is predicted to hold exactly for homogeneous ideal gases. Monatomic ideal gases (the six noble gases) possess three degrees of freedom per atom, corresponding to the three spatial directions. According to the equipartition of energy this means that there is a thermal energy of per atom. This corresponds very well with experimental data. The thermal energy can be used to calculate the root-mean-square speed of the atoms, which turns out to be inversely proportional to the square root of the atomic mass. The root mean square speeds found at room temperature accurately reflect this, ranging from for helium, down to for xenon.
Kinetic theory gives the average pressure for an ideal gas as
Combination with the ideal gas law
shows that the average translational kinetic energy is
Considering that the translational motion velocity vector has three degrees of freedom (one for each dimension) gives the average energy per degree of freedom equal to one third of that, i.e. .
The ideal gas equation is also obeyed closely by molecular gases; but the form for the heat capacity is more complicated, because the molecules possess additional internal degrees of freedom, as well as the three degrees of freedom for movement of the molecule as a whole. Diatomic gases, for example, possess a total of six degrees of simple freedom per molecule that are related to atomic motion (three translational, two rotational, and one vibrational). At lower temperatures, not all these degrees of freedom may fully participate in the gas heat capacity, due to quantum mechanical limits on the availability of excited states at the relevant thermal energy per molecule.
Role in Boltzmann factors
More generally, systems in equilibrium at temperature have probability of occupying a state with energy weighted by the corresponding Boltzmann factor:
where is the partition function. Again, it is the energy-like quantity that takes central importance.
Consequences of this include (in addition to the results for ideal gases above) the Arrhenius equation in chemical kinetics.
Role in the statistical definition of entropy
In statistical mechanics, the entropy of an isolated system at thermodynamic equilibrium is defined as the natural logarithm of , the number of distinct microscopic states available to the system given the macroscopic constraints (such as a fixed total energy ):
This equation, which relates the microscopic details, or microstates, of the system (via ) to its macroscopic state (via the entropy ), is the central idea of statistical mechanics. Such is its importance that it is inscribed on Boltzmann's tombstone.
The constant of proportionality serves to make the statistical mechanical entropy equal to the classical thermodynamic entropy of Clausius:
One could choose instead a rescaled dimensionless entropy in microscopic terms such that
This is a more natural form and this rescaled entropy exactly corresponds to Shannon's subsequent information entropy.
The characteristic energy is thus the energy required to increase the rescaled entropy by one nat.
Thermal voltage
In semiconductors, the Shockley diode equation—the relationship between the flow of electric current and the electrostatic potential across a p–n junction—depends on a characteristic voltage called the thermal voltage, denoted by . The thermal voltage depends on absolute temperature as
where is the magnitude of the electrical charge on the electron with a value Equivalently,
At room temperature , is approximately which can be derived by plugging in the values as follows:
At the standard state temperature of , it is approximately . The thermal voltage is also important in plasmas and electrolyte solutions (e.g. the Nernst equation); in both cases it provides a measure of how much the spatial distribution of electrons or ions is affected by a boundary held at a fixed voltage.
History
The Boltzmann constant is named after its 19th century Austrian discoverer, Ludwig Boltzmann. Although Boltzmann first linked entropy and probability in 1877, the relation was never expressed with a specific constant until Max Planck first introduced , and gave a more precise value for it (, about 2.5% lower than today's figure), in his derivation of the law of black-body radiation in 1900–1901. Before 1900, equations involving Boltzmann factors were not written using the energies per molecule and the Boltzmann constant, but rather using a form of the gas constant , and macroscopic energies for macroscopic quantities of the substance. The iconic terse form of the equation on Boltzmann's tombstone is in fact due to Planck, not Boltzmann. Planck actually introduced it in the same work as his eponymous .
In 1920, Planck wrote in his Nobel Prize lecture:
This "peculiar state of affairs" is illustrated by reference to one of the great scientific debates of the time. There was considerable disagreement in the second half of the nineteenth century as to whether atoms and molecules were real or whether they were simply a heuristic tool for solving problems. There was no agreement whether chemical molecules, as measured by atomic weights, were the same as physical molecules, as measured by kinetic theory. Planck's 1920 lecture continued:
In versions of SI prior to the 2019 revision of the SI, the Boltzmann constant was a measured quantity rather than having a fixed numerical value. Its exact definition also varied over the years due to redefinitions of the kelvin (see ) and other SI base units (see ).
In 2017, the most accurate measures of the Boltzmann constant were obtained by acoustic gas thermometry, which determines the speed of sound of a monatomic gas in a triaxial ellipsoid chamber using microwave and acoustic resonances. This decade-long effort was undertaken with different techniques by several laboratories; it is one of the cornerstones of the 2019 revision of the SI. Based on these measurements, the CODATA recommended to be the final fixed value of the Boltzmann constant to be used for the International System of Units.
As a precondition for redefining the Boltzmann constant, there must be one experimental value with a relative uncertainty below 1 ppm, and at least one measurement from a second technique with a relative uncertainty below 3 ppm. The acoustic gas thermometry reached 0.2 ppm, and Johnson noise thermometry reached 2.8 ppm.
Value in different units
Since is a proportionality factor between temperature and energy, its numerical value depends on the choice of units for energy and temperature. The small numerical value of the Boltzmann constant in SI units means a change in temperature by 1 K only changes a particle's energy by a small amount. A change of is defined to be the same as a change of . The characteristic energy is a term encountered in many physical relationships.
The Boltzmann constant sets up a relationship between wavelength and temperature (dividing hc/k by a wavelength gives a temperature) with one micrometer being related to , and also a relationship between voltage and temperature (kT in units of eV corresponds to a voltage) with one volt being related to . The ratio of these two temperatures, / ≈ 1.239842, is the numerical value of hc in units of eV⋅μm.
Natural units
The Boltzmann constant provides a mapping from the characteristic microscopic energy to the macroscopic temperature scale . In fundamental physics, this mapping is often simplified by using the natural units of setting to unity. This convention means that temperature and energy quantities have the same dimensions. In particular, the SI unit kelvin becomes superfluous, being defined in terms of joules as . With this convention, temperature is always given in units of energy, and the Boltzmann constant is not explicitly needed in formulas.
This convention simplifies many physical relationships and formulas. For example, the equipartition formula for the energy associated with each classical degree of freedom ( above) becomes
As another example, the definition of thermodynamic entropy coincides with the form of information entropy:
where is the probability of each microstate.
| Physical sciences | Physical constants | Physics |
53741 | https://en.wikipedia.org/wiki/Symmetry | Symmetry | Symmetry () in everyday life refers to a sense of harmonious and beautiful proportion and balance. In mathematics, the term has a more precise definition and is usually used to refer to an object that is invariant under some transformations, such as translation, reflection, rotation, or scaling. Although these two meanings of the word can sometimes be told apart, they are intricately related, and hence are discussed together in this article.
Mathematical symmetry may be observed with respect to the passage of time; as a spatial relationship; through geometric transformations; through other kinds of functional transformations; and as an aspect of abstract objects, including theoretic models, language, and music.
This article describes symmetry from three perspectives: in mathematics, including geometry, the most familiar type of symmetry for many people; in science and nature; and in the arts, covering architecture, art, and music.
The opposite of symmetry is asymmetry, which refers to the absence of symmetry.
In mathematics
In geometry
A geometric shape or object is symmetric if it can be divided into two or more identical pieces that are arranged in an organized fashion. This means that an object is symmetric if there is a transformation that moves individual pieces of the object, but doesn't change the overall shape. The type of symmetry is determined by the way the pieces are organized, or by the type of transformation:
An object has reflectional symmetry (line or mirror symmetry) if there is a line (or in 3D a plane) going through it which divides it into two pieces that are mirror images of each other.
An object has rotational symmetry if the object can be rotated about a fixed point (or in 3D about a line) without changing the overall shape.
An object has translational symmetry if it can be translated (moving every point of the object by the same distance) without changing its overall shape.
An object has helical symmetry if it can be simultaneously translated and rotated in three-dimensional space along a line known as a screw axis.
An object has scale symmetry if it does not change shape when it is expanded or contracted. Fractals also exhibit a form of scale symmetry, where smaller portions of the fractal are similar in shape to larger portions.
Other symmetries include glide reflection symmetry (a reflection followed by a translation) and rotoreflection symmetry (a combination of a rotation and a reflection).
In logic
A dyadic relation R = S × S is symmetric if for all elements a, b in S, whenever it is true that Rab, it is also true that Rba. Thus, the relation "is the same age as" is symmetric, for if Paul is the same age as Mary, then Mary is the same age as Paul.
In propositional logic, symmetric binary logical connectives include and (∧, or &), or (∨, or |) and if and only if (↔), while the connective if (→) is not symmetric. Other symmetric logical connectives include nand (not-and, or ⊼), xor (not-biconditional, or ⊻), and nor (not-or, or ⊽).
Other areas of mathematics
Generalizing from geometrical symmetry in the previous section, one can say that a mathematical object is symmetric with respect to a given mathematical operation, if, when applied to the object, this operation preserves some property of the object. The set of operations that preserve a given property of the object form a group.
In general, every kind of structure in mathematics will have its own kind of symmetry. Examples include even and odd functions in calculus, symmetric groups in abstract algebra, symmetric matrices in linear algebra, and Galois groups in Galois theory. In statistics, symmetry also manifests as symmetric probability distributions, and as skewness—the asymmetry of distributions.
In science and nature
In physics
Symmetry in physics has been generalized to mean invariance—that is, lack of change—under any kind of transformation, for example arbitrary coordinate transformations. This concept has become one of the most powerful tools of theoretical physics, as it has become evident that practically all laws of nature originate in symmetries. In fact, this role inspired the Nobel laureate PW Anderson to write in his widely read 1972 article More is Different that "it is only slightly overstating the case to say that physics is the study of symmetry." See Noether's theorem (which, in greatly simplified form, states that for every continuous mathematical symmetry, there is a corresponding conserved quantity such as energy or momentum; a conserved current, in Noether's original language); and also, Wigner's classification, which says that the symmetries of the laws of physics determine the properties of the particles found in nature.
Important symmetries in physics include continuous symmetries and discrete symmetries of spacetime; internal symmetries of particles; and supersymmetry of physical theories.
In biology
In biology, the notion of symmetry is mostly used explicitly to describe body shapes. Bilateral animals, including humans, are more or less symmetric with respect to the sagittal plane which divides the body into left and right halves. Animals that move in one direction necessarily have upper and lower sides, head and tail ends, and therefore a left and a right. The head becomes specialized with a mouth and sense organs, and the body becomes bilaterally symmetric for the purpose of movement, with symmetrical pairs of muscles and skeletal elements, though internal organs often remain asymmetric.
Plants and sessile (attached) animals such as sea anemones often have radial or rotational symmetry, which suits them because food or threats may arrive from any direction. Fivefold symmetry is found in the echinoderms, the group that includes starfish, sea urchins, and sea lilies.
In biology, the notion of symmetry is also used as in physics, that is to say to describe the properties of the objects studied, including their interactions. A remarkable property of biological evolution is the changes of symmetry corresponding to the appearance of new parts and dynamics.
In chemistry
Symmetry is important to chemistry because it undergirds essentially all specific interactions between molecules in nature (i.e., via the interaction of natural and human-made chiral molecules with inherently chiral biological systems). The control of the symmetry of molecules produced in modern chemical synthesis contributes to the ability of scientists to offer therapeutic interventions with minimal side effects. A rigorous understanding of symmetry explains fundamental observations in quantum chemistry, and in the applied areas of spectroscopy and crystallography. The theory and application of symmetry to these areas of physical science draws heavily on the mathematical area of group theory.
In psychology and neuroscience
For a human observer, some symmetry types are more salient than others, in particular the most salient is a reflection with a vertical axis, like that present in the human face. Ernst Mach made this observation in his book "The analysis of sensations" (1897), and this implies that perception of symmetry is not a general response to all types of regularities. Both behavioural and neurophysiological studies have confirmed the special sensitivity to reflection symmetry in humans and also in other animals. Early studies within the Gestalt tradition suggested that bilateral symmetry was one of the key factors in perceptual grouping. This is known as the Law of Symmetry. The role of symmetry in grouping and figure/ground organization has been confirmed in many studies. For instance, detection of reflectional symmetry is faster when this is a property of a single object. Studies of human perception and psychophysics have shown that detection of symmetry is fast, efficient and robust to perturbations. For example, symmetry can be detected with presentations between 100 and 150 milliseconds.
More recent neuroimaging studies have documented which brain regions are active during perception of symmetry. Sasaki et al. used functional magnetic resonance imaging (fMRI) to compare responses for patterns with symmetrical or random dots. A strong activity was present in extrastriate regions of the occipital cortex but not in the primary visual cortex. The extrastriate regions included V3A, V4, V7, and the lateral occipital complex (LOC). Electrophysiological studies have found a late posterior negativity that originates from the same areas. In general, a large part of the visual system seems to be involved in processing visual symmetry, and these areas involve similar networks to those responsible for detecting and recognising objects.
In social interactions
People observe the symmetrical nature, often including asymmetrical balance, of social interactions in a variety of contexts. These include assessments of reciprocity, empathy, sympathy, apology, dialogue, respect, justice, and revenge.
Reflective equilibrium is the balance that may be attained through deliberative mutual adjustment among general principles and specific judgments.
Symmetrical interactions send the moral message "we are all the same" while asymmetrical interactions may send the message "I am special; better than you." Peer relationships, such as can be governed by the Golden Rule, are based on symmetry, whereas power relationships are based on asymmetry. Symmetrical relationships can to some degree be maintained by simple (game theory) strategies seen in symmetric games such as tit for tat.
In the arts
There exists a list of journals and newsletters known to deal, at least in part, with symmetry and the arts.
In architecture
Symmetry finds its ways into architecture at every scale, from the overall external views of buildings such as Gothic cathedrals and The White House, through the layout of the individual floor plans, and down to the design of individual building elements such as tile mosaics. Islamic buildings such as the Taj Mahal and the Lotfollah mosque make elaborate use of symmetry both in their structure and in their ornamentation. Moorish buildings like the Alhambra are ornamented with complex patterns made using translational and reflection symmetries as well as rotations.
It has been said that only bad architects rely on a "symmetrical layout of blocks, masses and structures"; Modernist architecture, starting with International style, relies instead on "wings and balance of masses".
In pottery and metal vessels
Since the earliest uses of pottery wheels to help shape clay vessels, pottery has had a strong relationship to symmetry. Pottery created using a wheel acquires full rotational symmetry in its cross-section, while allowing substantial freedom of shape in the vertical direction. Upon this inherently symmetrical starting point, potters from ancient times onwards have added patterns that modify the rotational symmetry to achieve visual objectives.
Cast metal vessels lacked the inherent rotational symmetry of wheel-made pottery, but otherwise provided a similar opportunity to decorate their surfaces with patterns pleasing to those who used them. The ancient Chinese, for example, used symmetrical patterns in their bronze castings as early as the 17th century BC. Bronze vessels exhibited both a bilateral main motif and a repetitive translated border design.
In carpets and rugs
A long tradition of the use of symmetry in carpet and rug patterns spans a variety of cultures. American Navajo Indians used bold diagonals and rectangular motifs. Many Oriental rugs have intricate reflected centers and borders that translate a pattern. Not surprisingly, rectangular rugs have typically the symmetries of a rectangle—that is, motifs that are reflected across both the horizontal and vertical axes (see ).
In quilts
As quilts are made from square blocks (usually 9, 16, or 25 pieces to a block) with each smaller piece usually consisting of fabric triangles, the craft lends itself readily to the application of symmetry.
In other arts and crafts
Symmetries appear in the design of objects of all kinds. Examples include beadwork, furniture, sand paintings, knotwork, masks, and musical instruments. Symmetries are central to the art of M.C. Escher and the many applications of tessellation in art and craft forms such as wallpaper, ceramic tilework such as in Islamic geometric decoration, batik, ikat, carpet-making, and many kinds of textile and embroidery patterns.
Symmetry is also used in designing logos. By creating a logo on a grid and using the theory of symmetry, designers can organize their work, create a symmetric or asymmetrical design, determine the space between letters, determine how much negative space is required in the design, and how to accentuate parts of the logo to make it stand out.
In music
Symmetry is not restricted to the visual arts. Its role in the history of music touches many aspects of the creation and perception of music.
Musical form
Symmetry has been used as a formal constraint by many composers, such as the arch (swell) form (ABCBA) used by Steve Reich, Béla Bartók, and James Tenney. In classical music, Johann Sebastian Bach used the symmetry concepts of permutation and invariance.
Pitch structures
Symmetry is also an important consideration in the formation of scales and chords, traditional or tonal music being made up of non-symmetrical groups of pitches, such as the diatonic scale or the major chord. Symmetrical scales or chords, such as the whole tone scale, augmented chord, or diminished seventh chord (diminished-diminished seventh), are said to lack direction or a sense of forward motion, are ambiguous as to the key or tonal center, and have a less specific diatonic functionality. However, composers such as Alban Berg, Béla Bartók, and George Perle have used axes of symmetry and/or interval cycles in an analogous way to keys or non-tonal tonal centers. George Perle explains that "C–E, D–F♯, [and] Eb–G, are different instances of the same interval … the other kind of identity. … has to do with axes of symmetry. C–E belongs to a family of symmetrically related dyads as follows:"
Thus in addition to being part of the interval-4 family, C–E is also a part of the sum-4 family (with C equal to 0).
Interval cycles are symmetrical and thus non-diatonic. However, a seven pitch segment of C5 (the cycle of fifths, which are enharmonic with the cycle of fourths) will produce the diatonic major scale. Cyclic tonal progressions in the works of Romantic composers such as Gustav Mahler and Richard Wagner form a link with the cyclic pitch successions in the atonal music of Modernists such as Bartók, Alexander Scriabin, Edgard Varèse, and the Vienna school. At the same time, these progressions signal the end of tonality.
The first extended composition consistently based on symmetrical pitch relations was probably Alban Berg's Quartet, Op. 3 (1910).
Equivalency
Tone rows or pitch class sets which are invariant under retrograde are horizontally symmetrical, under inversion vertically. | Mathematics | Geometry | null |
53759 | https://en.wikipedia.org/wiki/Category%20%28mathematics%29 | Category (mathematics) | In mathematics, a category (sometimes called an abstract category to distinguish it from a concrete category) is a collection of "objects" that are linked by "arrows". A category has two basic properties: the ability to compose the arrows associatively and the existence of an identity arrow for each object. A simple example is the category of sets, whose objects are sets and whose arrows are functions.
Category theory is a branch of mathematics that seeks to generalize all of mathematics in terms of categories, independent of what their objects and arrows represent. Virtually every branch of modern mathematics can be described in terms of categories, and doing so often reveals deep insights and similarities between seemingly different areas of mathematics. As such, category theory provides an alternative foundation for mathematics to set theory and other proposed axiomatic foundations. In general, the objects and arrows may be abstract entities of any kind, and the notion of category provides a fundamental and abstract way to describe mathematical entities and their relationships.
In addition to formalizing mathematics, category theory is also used to formalize many other systems in computer science, such as the semantics of programming languages.
Two categories are the same if they have the same collection of objects, the same collection of arrows, and the same associative method of composing any pair of arrows. Two different categories may also be considered "equivalent" for purposes of category theory, even if they do not have precisely the same structure.
Well-known categories are denoted by a short capitalized word or abbreviation in bold or italics: examples include Set, the category of sets and set functions; Ring, the category of rings and ring homomorphisms; and Top, the category of topological spaces and continuous maps. All of the preceding categories have the identity map as identity arrows and composition as the associative operation on arrows.
The classic and still much used text on category theory is Categories for the Working Mathematician by Saunders Mac Lane. Other references are given in the | Mathematics | Category theory | null |
53776 | https://en.wikipedia.org/wiki/Kruskal%27s%20algorithm | Kruskal's algorithm | Kruskal's algorithm finds a minimum spanning forest of an undirected edge-weighted graph. If the graph is connected, it finds a minimum spanning tree. It is a greedy algorithm that in each step adds to the forest the lowest-weight edge that will not form a cycle. The key steps of the algorithm are sorting and the use of a disjoint-set data structure to detect cycles. Its running time is dominated by the time to sort all of the graph edges by their weight.
A minimum spanning tree of a connected weighted graph is a connected subgraph, without cycles, for which the sum of the weights of all the edges in the subgraph is minimal. For a disconnected graph, a minimum spanning forest is composed of a minimum spanning tree for each connected component.
This algorithm was first published by Joseph Kruskal in 1956, and was rediscovered soon afterward by . Other algorithms for this problem include Prim's algorithm, Borůvka's algorithm, and the reverse-delete algorithm.
Algorithm
The algorithm performs the following steps:
Create a forest (a set of trees) initially consisting of a separate single-vertex tree for each vertex in the input graph.
Sort the graph edges by weight.
Loop through the edges of the graph, in ascending sorted order by their weight. For each edge:
Test whether adding the edge to the current forest would create a cycle.
If not, add the edge to the forest, combining two trees into a single tree.
At the termination of the algorithm, the forest forms a minimum spanning forest of the graph. If the graph is connected, the forest has a single component and forms a minimum spanning tree.
Pseudocode
The following code is implemented with a disjoint-set data structure. It represents the forest F as a set of undirected edges, and uses the disjoint-set data structure to efficiently determine whether two vertices are part of the same tree.
algorithm Kruskal(G) is
F:= ∅
for each v in G.V do
MAKE-SET(v)
for each {u, v} in G.E ordered by weight({u, v}), increasing do
if FIND-SET(u) ≠ FIND-SET(v) then
F := F ∪ { {u, v} }
UNION(FIND-SET(u), FIND-SET(v))
return F
Complexity
For a graph with edges and vertices, Kruskal's algorithm can be shown to run in time time, with simple data structures. Here, expresses the time in big O notation, and is a logarithm to any base (since inside -notation logarithms to all bases are equivalent, because they are the same up to a constant factor). This time bound is often written instead as , which is equivalent for graphs with no isolated vertices, because for these graphs and the logarithms of and are again within a constant factor of each other.
To achieve this bound, first sort the edges by weight using a comparison sort in time. Once sorted, it is possible to loop through the edges in sorted order in constant time per edge. Next, use a disjoint-set data structure, with a set of vertices for each component, to keep track of which vertices are in which components. Creating this structure, with a separate set for each vertex, takes operations and time. The final iteration through all edges performs two find operations and possibly one union operation per edge. These operations take amortized time time per operation, giving worst-case total time for this loop, where is the extremely slowly growing inverse Ackermann function. This part of the time bound is much smaller than the time for the sorting step, so the total time for the algorithm can be simplified to the time for the sorting step.
In cases where the edges are already sorted, or where they have small enough integer weight to allow integer sorting algorithms such as counting sort or radix sort to sort them in linear time, the disjoint set operations are the slowest remaining part of the algorithm and the total time is {{math|O(E α(V))}}.
Example
Proof of correctness
The proof consists of two parts. First, it is proved that the algorithm produces a spanning tree. Second, it is proved that the constructed spanning tree is of minimal weight.
Spanning tree
Let be a connected, weighted graph and let be the subgraph of produced by the algorithm. cannot have a cycle, as by definition an edge is not added if it results in a cycle. cannot be disconnected, since the first encountered edge that joins two components of would have been added by the algorithm. Thus, is a spanning tree of .
Minimality
We show that the following proposition P is true by induction: If F is the set of edges chosen at any stage of the algorithm, then there is some minimum spanning tree that contains F and none of the edges rejected by the algorithm.
Clearly P is true at the beginning, when F is empty: any minimum spanning tree will do, and there exists one because a weighted connected graph always has a minimum spanning tree.
Now assume P is true for some non-final edge set F and let T be a minimum spanning tree that contains F.
If the next chosen edge e is also in T, then P is true for F + e.
Otherwise, if e is not in T then T + e has a cycle C. The cycle C contains edges which do not belong to F + e, since e does not form a cycle when added to F but does in T. Let f be an edge which is in C but not in F + e. Note that f also belongs to T, since f belongs to T + e but not F + e. By P, f has not been considered by the algorithm. f must therefore have a weight at least as large as e. Then T − f + e is a tree, and it has the same or less weight as T. However since T is a minimum spanning tree then T − f + e has the same weight as T, otherwise we get a contradiction and T would not be a minimum spanning tree. So T − f + e is a minimum spanning tree containing F + e and again P holds.
Therefore, by the principle of induction, P holds when F has become a spanning tree, which is only possible if F is a minimum spanning tree itself.
Parallel algorithm
Kruskal's algorithm is inherently sequential and hard to parallelize. It is, however, possible to perform the initial sorting of the edges in parallel or, alternatively, to use a parallel implementation of a binary heap to extract the minimum-weight edge in every iteration.
As parallel sorting is possible in time on processors, the runtime of Kruskal's algorithm can be reduced to O(E α(V)), where α again is the inverse of the single-valued Ackermann function.
A variant of Kruskal's algorithm, named Filter-Kruskal, has been described by Osipov et al. and is better suited for parallelization. The basic idea behind Filter-Kruskal is to partition the edges in a similar way to quicksort and filter out edges that connect vertices of the same tree to reduce the cost of sorting. The following pseudocode demonstrates this.
function filter_kruskal(G) is
if |G.E| < kruskal_threshold:
return kruskal(G)
pivot = choose_random(G.E)
E, E = partition(G.E, pivot)
A = filter_kruskal(E)
E = filter(E)
A = A ∪ filter_kruskal(E)
return A
function partition(E, pivot) is
E = ∅, E = ∅
foreach (u, v) in E do
if weight(u, v) ≤ pivot then
E = E ∪ {(u, v)}
else
E = E ∪ {(u, v)}
return E, E
function filter(E) is
E = ∅
foreach (u, v) in E do
if find_set(u) ≠ find_set(v) then
E = E ∪ {(u, v)}
return E
Filter-Kruskal lends itself better to parallelization as sorting, filtering, and partitioning can easily be performed in parallel by distributing the edges between the processors.
Finally, other variants of a parallel implementation of Kruskal's algorithm have been explored. Examples include a scheme that uses helper threads to remove edges that are definitely not part of the MST in the background, and a variant which runs the sequential algorithm on p subgraphs, then merges those subgraphs until only one, the final MST, remains.
| Mathematics | Graph theory | null |
53777 | https://en.wikipedia.org/wiki/Ford%E2%80%93Fulkerson%20algorithm | Ford–Fulkerson algorithm | The Ford–Fulkerson method or Ford–Fulkerson algorithm (FFA) is a greedy algorithm that computes the maximum flow in a flow network. It is sometimes called a "method" instead of an "algorithm" as the approach to finding augmenting paths in a residual graph is not fully specified or it is specified in several implementations with different running times. It was published in 1956 by L. R. Ford Jr. and D. R. Fulkerson. The name "Ford–Fulkerson" is often also used for the Edmonds–Karp algorithm, which is a fully defined implementation of the Ford–Fulkerson method.
The idea behind the algorithm is as follows: as long as there is a path from the source (start node) to the sink (end node), with available capacity on all edges in the path, we send flow along one of the paths. Then we find another path, and so on. A path with available capacity is called an augmenting path.
Algorithm
Let be a graph, and for each edge from to , let be the capacity and be the flow. We want to find the maximum flow from the source to the sink . After every step in the algorithm the following is maintained:
{| class="wikitable"
! | Capacity constraints
| || The flow along an edge cannot exceed its capacity.
|-
! | Skew symmetry
| || The net flow from to must be the opposite of the net flow from to (see example).
|-
! | Flow conservation
| || The net flow to a node is zero, except for the source, which "produces" flow, and the sink, which "consumes" flow.
|-
! | Value(f)
| || The flow leaving from must be equal to the flow arriving at .
|-
|}
This means that the flow through the network is a legal flow after each round in the algorithm. We define the residual network to be the network with capacity and no flow. Notice that it can happen that a flow from to is allowed in the residual
network, though disallowed in the original network: if and then .
Inputs Given a Network with flow capacity , a source node , and a sink node
Output Compute a flow from to of maximum value
for all edges
While there is a path from to in , such that for all edges :
Find
For each edge
(Send flow along the path)
(The flow might be "returned" later)
The path in step 2 can be found with, for example, breadth-first search (BFS) or depth-first search in . The former is known as the Edmonds–Karp algorithm.
When no more paths in step 2 can be found, will not be able to reach in the residual
network. If is the set of nodes reachable by in the residual network, then the total
capacity in the original network of edges from to the remainder of is on the one hand
equal to the total flow we found from to ,
and on the other hand serves as an upper bound for all such flows.
This proves that the flow we found is maximal. | Mathematics | Graph theory | null |
53783 | https://en.wikipedia.org/wiki/Prim%27s%20algorithm | Prim's algorithm | In computer science, Prim's algorithm is a greedy algorithm that finds a minimum spanning tree for a weighted undirected graph. This means it finds a subset of the edges that forms a tree that includes every vertex, where the total weight of all the edges in the tree is minimized. The algorithm operates by building this tree one vertex at a time, from an arbitrary starting vertex, at each step adding the cheapest possible connection from the tree to another vertex.
The algorithm was developed in 1930 by Czech mathematician Vojtěch Jarník and later rediscovered and republished by computer scientists Robert C. Prim in 1957 and Edsger W. Dijkstra in 1959. Therefore, it is also sometimes called the Jarník's algorithm, Prim–Jarník algorithm, Prim–Dijkstra algorithm
or the DJP algorithm.
Other well-known algorithms for this problem include Kruskal's algorithm and Borůvka's algorithm. These algorithms find the minimum spanning forest in a possibly disconnected graph; in contrast, the most basic form of Prim's algorithm only finds minimum spanning trees in connected graphs. However, running Prim's algorithm separately for each connected component of the graph, it can also be used to find the minimum spanning forest. In terms of their asymptotic time complexity, these three algorithms are equally fast for sparse graphs, but slower than other more sophisticated algorithms.
However, for graphs that are sufficiently dense, Prim's algorithm can be made to run in linear time, meeting or improving the time bounds for other algorithms.
Description
The algorithm may informally be described as performing the following steps:
In more detail, it may be implemented following the pseudocode below.
As described above, the starting vertex for the algorithm will be chosen arbitrarily, because the first iteration of the main loop of the algorithm will have a set of vertices in Q that all have equal weights, and the algorithm will automatically start a new tree in F when it completes a spanning tree of each connected component of the input graph. The algorithm may be modified to start with any particular vertex s by setting C[s] to be a number smaller than the other values of C (for instance, zero), and it may be modified to only find a single spanning tree rather than an entire spanning forest (matching more closely the informal description) by stopping whenever it encounters another vertex flagged as having no associated edge.
Different variations of the algorithm differ from each other in how the set Q is implemented: as a simple linked list or array of vertices, or as a more complicated priority queue data structure. This choice leads to differences in the time complexity of the algorithm. In general, a priority queue will be quicker at finding the vertex v with minimum cost, but will entail more expensive updates when the value of C[w] changes.
Time complexity
The time complexity of Prim's algorithm depends on the data structures used for the graph and for ordering the edges by weight, which can be done using a priority queue. The following table shows the typical choices:
A simple implementation of Prim's, using an adjacency matrix or an adjacency list graph representation and linearly searching an array of weights to find the minimum weight edge to add, requires O(|V|2) running time. However, this running time can be greatly improved by using heaps to implement finding minimum weight edges in the algorithm's inner loop.
A first improved version uses a heap to store all edges of the input graph, ordered by their weight. This leads to an O(|E| log |E|) worst-case running time. But storing vertices instead of edges can improve it still further. The heap should order the vertices by the smallest edge-weight that connects them to any vertex in the partially constructed minimum spanning tree (MST) (or infinity if no such edge exists). Every time a vertex v is chosen and added to the MST, a decrease-key operation is performed on all vertices w outside the partial MST such that v is connected to w, setting the key to the minimum of its previous value and the edge cost of (v,w).
Using a simple binary heap data structure, Prim's algorithm can now be shown to run in time O(|E| log |V|) where |E| is the number of edges and |V| is the number of vertices. Using a more sophisticated Fibonacci heap, this can be brought down to O(|E| + |V| log |V|), which is asymptotically faster when the graph is dense enough that |E| is ω(|V|), and linear time when |E| is at least |V| log |V|. For graphs of even greater density (having at least |V|c edges for some c > 1), Prim's algorithm can be made to run in linear time even more simply, by using a d-ary heap in place of a Fibonacci heap.
Proof of correctness
Let P be a connected, weighted graph. At every iteration of Prim's algorithm, an edge must be found that connects a vertex in a subgraph to a vertex outside the subgraph. Since P is connected, there will always be a path to every vertex. The output Y of Prim's algorithm is a tree, because the edge and vertex added to tree Y are connected. Let Y1 be a minimum spanning tree of graph P. If Y1=Y then Y is a minimum spanning tree. Otherwise, let e be the first edge added during the construction of tree Y that is not in tree Y1, and V be the set of vertices connected by the edges added before edge e. Then one endpoint of edge e is in set V and the other is not. Since tree Y1 is a spanning tree of graph P, there is a path in tree Y1 joining the two endpoints. As one travels along the path, one must encounter an edge f joining a vertex in set V to one that is not in set V. Now, at the iteration when edge e was added to tree Y, edge f could also have been added and it would be added instead of edge e if its weight was less than e, and since edge f was not added, we conclude that
Let tree Y2 be the graph obtained by removing edge f from and adding edge e to tree Y1. It is easy to show that tree Y2 is connected, has the same number of edges as tree Y1, and the total weights of its edges is not larger than that of tree Y1, therefore it is also a minimum spanning tree of graph P and it contains edge e and all the edges added before it during the construction of set V. Repeat the steps above and we will eventually obtain a minimum spanning tree of graph P that is identical to tree Y. This shows Y is a minimum spanning tree. The minimum spanning tree allows for the first subset of the sub-region to be expanded into a larger subset X, which we assume to be the minimum.
Parallel algorithm
The main loop of Prim's algorithm is inherently sequential and thus not parallelizable. However, the inner loop, which determines the next edge of minimum weight that does not form a cycle, can be parallelized by dividing the vertices and edges between the available processors. The following pseudocode demonstrates this.
This algorithm can generally be implemented on distributed machines as well as on shared memory machines. The running time is , assuming that the reduce and broadcast operations can be performed in . A variant of Prim's algorithm for shared memory machines, in which Prim's sequential algorithm is being run in parallel, starting from different vertices, has also been explored. It should, however, be noted that more sophisticated algorithms exist to solve the distributed minimum spanning tree problem in a more efficient manner.
| Mathematics | Graph theory | null |
53784 | https://en.wikipedia.org/wiki/Brute-force%20attack | Brute-force attack | In cryptography, a brute-force attack consists of an attacker submitting many passwords or passphrases with the hope of eventually guessing correctly. The attacker systematically checks all possible passwords and passphrases until the correct one is found. Alternatively, the attacker can attempt to guess the key which is typically created from the password using a key derivation function. This is known as an exhaustive key search. This approach doesn't depend on intellectual tactics; rather, it relies on making several attempts.
A brute-force attack is a cryptanalytic attack that can, in theory, be used to attempt to decrypt any encrypted data (except for data encrypted in an information-theoretically secure manner). Such an attack might be used when it is not possible to take advantage of other weaknesses in an encryption system (if any exist) that would make the task easier.
When password-guessing, this method is very fast when used to check all short passwords, but for longer passwords other methods such as the dictionary attack are used because a brute-force search takes too long. Longer passwords, passphrases and keys have more possible values, making them exponentially more difficult to crack than shorter ones due to diversity of characters.
Brute-force attacks can be made less effective by obfuscating the data to be encoded making it more difficult for an attacker to recognize when the code has been cracked or by making the attacker do more work to test each guess. One of the measures of the strength of an encryption system is how long it would theoretically take an attacker to mount a successful brute-force attack against it.
Brute-force attacks are an application of brute-force search, the general problem-solving technique of enumerating all candidates and checking each one. The word 'hammering' is sometimes used to describe a brute-force attack, with 'anti-hammering' for countermeasures.
Basic concept
Brute-force attacks work by calculating every possible combination that could make up a password and testing it to see if it is the correct password. As the password's length increases, the amount of time, on average, to find the correct password increases exponentially.
Theoretical limits
The resources required for a brute-force attack grow exponentially with increasing key size, not linearly. Although U.S. export regulations historically restricted key lengths to 56-bit symmetric keys (e.g. Data Encryption Standard), these restrictions are no longer in place, so modern symmetric algorithms typically use computationally stronger 128- to 256-bit keys.
There is a physical argument that a 128-bit symmetric key is computationally secure against brute-force attack. The Landauer limit implied by the laws of physics sets a lower limit on the energy required to perform a computation of per bit erased in a computation, where T is the temperature of the computing device in kelvins, k is the Boltzmann constant, and the natural logarithm of 2 is about 0.693 (0.6931471805599453). No irreversible computing device can use less energy than this, even in principle. Thus, in order to simply flip through the possible values for a 128-bit symmetric key (ignoring doing the actual computing to check it) would, theoretically, require 2128 − 1 bit flips on a conventional processor. If it is assumed that the calculation occurs near room temperature (≈300 K), the Von Neumann-Landauer Limit can be applied to estimate the energy required as ≈1018 joules, which is equivalent to consuming 30 gigawatts of power for one year. This is equal to 30×109 W×365×24×3600 s = 9.46×1017 J or 262.7 TWh (about 0.1% of the yearly world energy production). The full actual computation – checking each key to see if a solution has been found – would consume many times this amount. Furthermore, this is simply the energy requirement for cycling through the key space; the actual time it takes to flip each bit is not considered, which is certainly greater than 0 (see Bremermann's limit).
However, this argument assumes that the register values are changed using conventional set and clear operations, which inevitably generate entropy. It has been shown that computational hardware can be designed not to encounter this theoretical obstruction (see reversible computing), though no such computers are known to have been constructed.
As commercial successors of governmental ASIC solutions have become available, also known as custom hardware attacks, two emerging technologies have proven their capability in the brute-force attack of certain ciphers. One is modern graphics processing unit (GPU) technology, the other is the field-programmable gate array (FPGA) technology. GPUs benefit from their wide availability and price-performance benefit, FPGAs from their energy efficiency per cryptographic operation. Both technologies try to transport the benefits of parallel processing to brute-force attacks. In case of GPUs some hundreds, in the case of FPGA some thousand processing units making them much better suited to cracking passwords than conventional processors. For instance in 2022, 8 Nvidia RTX 4090 GPU were linked together to test password strength by using the software Hashcat with results that showed 200 billion eight-character NTLM password combinations could be cycled through in 48 minutes.
Various publications in the fields of cryptographic analysis have proved the energy efficiency of today's FPGA technology, for example, the COPACOBANA FPGA Cluster computer consumes the same energy as a single PC (600 W), but performs like 2,500 PCs for certain algorithms. A number of firms provide hardware-based FPGA cryptographic analysis solutions from a single FPGA PCI Express card up to dedicated FPGA computers. WPA and WPA2 encryption have successfully been brute-force attacked by reducing the workload by a factor of 50 in comparison to conventional CPUs and some hundred in case of FPGAs.
Advanced Encryption Standard (AES) permits the use of 256-bit keys. Breaking a symmetric 256-bit key by brute-force requires 2128 times more computational power than a 128-bit key. One of the fastest supercomputers in 2019 has a speed of 100 petaFLOPS which could theoretically check 100 trillion (1014) AES keys per second (assuming 1000 operations per check), but would still require 3.67×1055 years to exhaust the 256-bit key space.
An underlying assumption of a brute-force attack is that the complete key space was used to generate keys, something that relies on an effective random number generator, and that there are no defects in the algorithm or its implementation. For example, a number of systems that were originally thought to be impossible to crack by brute-force have nevertheless been cracked because the key space to search through was found to be much smaller than originally thought, because of a lack of entropy in their pseudorandom number generators. These include Netscape's implementation of Secure Sockets Layer (SSL) (cracked by Ian Goldberg and David Wagner in 1995) and a Debian/Ubuntu edition of OpenSSL discovered in 2008 to be flawed. A similar lack of implemented entropy led to the breaking of Enigma's code.
Credential recycling
Credential recycling is the hacking practice of re-using username and password combinations gathered in previous brute-force attacks. A special form of credential recycling is pass the hash, where unsalted hashed credentials are stolen and re-used without first being brute-forced.
Unbreakable codes
Certain types of encryption, by their mathematical properties, cannot be defeated by brute-force. An example of this is one-time pad cryptography, where every cleartext bit has a corresponding key from a truly random sequence of key bits. A 140 character one-time-pad-encoded string subjected to a brute-force attack would eventually reveal every 140 character string possible, including the correct answer – but of all the answers given, there would be no way of knowing which was the correct one. Defeating such a system, as was done by the Venona project, generally relies not on pure cryptography, but upon mistakes in its implementation, such as the key pads not being truly random, intercepted keypads, or operators making mistakes.
Countermeasures
In case of an offline attack where the attacker has gained access to the encrypted material, one can try key combinations without the risk of discovery or interference. In case of online attacks, database and directory administrators can deploy countermeasures such as limiting the number of attempts that a password can be tried, introducing time delays between successive attempts, increasing the answer's complexity (e.g., requiring a CAPTCHA answer or employing multi-factor authentication), and/or locking accounts out after unsuccessful login attempts. Website administrators may prevent a particular IP address from trying more than a predetermined number of password attempts against any account on the site. Additionally, the MITRE D3FEND framework provides structured recommendations for defending against brute-force attacks by implementing strategies such as network traffic filtering, deploying decoy credentials, and invalidating authentication caches.
Reverse brute-force attack
In a reverse brute-force attack, a single (usually common) password is tested against multiple usernames or encrypted files. The process may be repeated for a select few passwords. In such a strategy, the attacker is not targeting a specific user.
| Technology | Computer security | null |
53878 | https://en.wikipedia.org/wiki/Betelgeuse | Betelgeuse | Betelgeuse is a red supergiant star in the constellation of Orion. It is usually the tenth-brightest star in the night sky and, after Rigel, the second-brightest in its constellation. It is a distinctly reddish, semiregular variable star whose apparent magnitude, varying between +0.0 and +1.6, has the widest range displayed by any first-magnitude star. Betelgeuse is the brightest star in the night sky at near-infrared wavelengths. Its Bayer designation is α Orionis, Latinised to Alpha Orionis and abbreviated Alpha Ori or α Ori.
With a radius between 640 and 764 times that of the Sun, if it were at the center of our Solar System, its surface would lie beyond the asteroid belt and it would engulf the orbits of Mercury, Venus, Earth, and Mars. Calculations of Betelgeuse's mass range from slightly under ten to a little over twenty times that of the Sun. For various reasons, its distance has been quite difficult to measure; current best estimates are of the order of 400–600 light-years from the Suna comparatively wide uncertainty for a relatively nearby star. Its absolute magnitude is about −6. With an age of less than 10 million years, Betelgeuse has evolved rapidly because of its large mass, and is expected to end its evolution with a supernova explosion, most likely within 100,000 years. When Betelgeuse explodes, it will shine as bright as the half-Moon for more than three months; life on Earth will be unharmed. Having been ejected from its birthplace in the Orion OB1 associationwhich includes the stars in Orion's Beltthis runaway star has been observed to be moving through the interstellar medium at a speed of , creating a bow shock over four light-years wide.
Betelgeuse became the first extrasolar star whose photosphere's angular size was measured in 1920, and subsequent studies have reported an angular diameter (i.e., apparent size) ranging from 0.042 to 0.056 arcseconds; that range of determinations is ascribed to non-sphericity, limb darkening, pulsations and varying appearance at different wavelengths. It is also surrounded by a complex, asymmetric envelope, roughly 250 times the size of the star, caused by mass loss from the star itself. The Earth-observed angular diameter of Betelgeuse is exceeded only by those of R Doradus and the Sun.
Starting in October 2019, Betelgeuse began to dim noticeably, and by mid-February 2020 its brightness had dropped by a factor of approximately 3, from magnitude 0.5 to 1.7. It then returned to a more normal brightness range, reaching a peak of 0.0 visual and 0.1 V-band magnitude in April 2023. Infrared observations found no significant change in luminosity over the last 50 years, suggesting that the dimming was due to a change in extinction around the star rather than a more fundamental change. A study using the Hubble Space Telescope suggests that occluding dust was created by a surface mass ejection; this material was cast millions of miles from the star, and then cooled to form the dust that caused the dimming.
Nomenclature
The star's designation is α Orionis (Latinised to Alpha Orionis), given by Johann Bayer in 1603.
The traditional name Betelgeuse was derived from the Arabic "the hand of al-Jawzā’ [i.e. Orion]". An error in the 13th-century reading of the Arabic initial yā’ () as bā’ (—a difference in i‘jām) led to the European name. In English, there are four common pronunciations of this name, depending on whether the first e is pronounced short or long and whether the s is pronounced or :
;
;
;
, popularized for sounding like "beetle juice".
In 2016, the International Astronomical Union organized a Working Group on Star Names (WGSN) to catalog and standardize proper names for stars. The WGSN's first bulletin, issued July 2016, included a table of the first two batches of names approved by the WGSN, which included Betelgeuse for this star. It is now so entered in the IAU Catalog of Star Names.
Observational history
Betelgeuse and its red coloration have been noted since antiquity; the classical astronomer Ptolemy described its color as (hypókirrhos = more or less orange-tawny), a term later described by a translator of Ulugh Beg's Zij-i Sultani as rubedo, Latin for "ruddiness". In the 19th century, before modern systems of stellar classification, Angelo Secchi included Betelgeuse as one of the prototypes for his Class III (orange to red) stars. Three centuries before Ptolemy, in contrast, Chinese astronomers observed Betelgeuse as yellow; Such an observation, if accurate, could suggest the star was in a yellow supergiant phase around this time, a credible possibility, given current research into these stars' complex circumstellar environment.
Nascent discoveries
Aboriginal groups in South Australia have shared oral tales of the variable brightness of Betelgeuse for an unknown period.
The variation in Betelgeuse's brightness was described in 1836 by Sir John Herschel in Outlines of Astronomy. From 1836 to 1840, he noticed significant changes in magnitude when Betelgeuse outshone Rigel in October 1837 and again in November 1839. A 10-year quiescent period followed; then in 1849, Herschel noted another short cycle of variability, which peaked in 1852. Later observers recorded unusually high maxima with an interval of years, but only small variations from 1957 to 1967. The records of the American Association of Variable Star Observers (AAVSO) show a maximum brightness of 0.2 in 1933 and 1942, and a minimum of 1.2, observed in 1927 and 1941. This variability in brightness may explain why Johann Bayer, with the publication of his Uranometria in 1603, designated the star alpha, as it probably rivaled the usually brighter Rigel (beta). From Arctic latitudes, Betelgeuse's red colour and higher location in the sky than Rigel meant the Inuit regarded it as brighter, and one local name was Ulluriajjuaq ("large star").
In 1920, Albert A. Michelson and Francis G. Pease mounted a six-meter interferometer on the front of the 2.5-meter telescope at Mount Wilson Observatory, helped by John August Anderson. The trio measured the angular diameter of Betelgeuse at 0.047″, a figure that resulted in a diameter of () based on the parallax value of . But limb darkening and measurement errors resulted in uncertainty about the accuracy of these measurements.
The 1950s and 1960s saw two developments that affected stellar convection theory in red supergiants: the Stratoscope projects and the 1958 publication of Structure and Evolution of the Stars, principally the work of Martin Schwarzschild and his colleague at Princeton University, Richard Härm.
This book disseminated ideas on how to apply computer technologies to create stellar models, while the Stratoscope projects, by taking balloon-borne telescopes above the Earth's turbulence, produced some of the finest images of solar granules and sunspots ever seen, thus confirming the existence of convection in the solar atmosphere.
Imaging breakthroughs
Astronomers saw some major advances in astronomical imaging technology in the 1970s, beginning with Antoine Labeyrie's invention of speckle interferometry, a process that significantly reduced the blurring effect caused by astronomical seeing. It increased the optical resolution of ground-based telescopes, allowing for more precise measurements of Betelgeuse's photosphere. With improvements in infrared telescopy atop Mount Wilson, Mount Locke, and Mauna Kea in Hawaii, astrophysicists began peering into the complex circumstellar shells surrounding the supergiant, causing them to suspect the presence of huge gas bubbles resulting from convection. However, it was not until the late 1980s and early 1990s, when Betelgeuse became a regular target for aperture masking interferometry, that breakthroughs occurred in visible-light and infrared imaging. Pioneered by J.E. Baldwin and colleagues of the Cavendish Astrophysics Group, the new technique employed a small mask with several holes in the telescope pupil plane, converting the aperture into an ad hoc interferometric array. The technique contributed some of the most accurate measurements of Betelgeuse while revealing bright spots on the star's photosphere. These were the first optical and infrared images of a stellar disk other than the Sun, taken first from ground-based interferometers and later from higher-resolution observations of the COAST telescope. The "bright patches" or "hotspots" observed with these instruments appeared to corroborate a theory put forth by Schwarzschild decades earlier of massive convection cells dominating the stellar surface.
In 1995, the Hubble Space Telescope's Faint Object Camera captured an ultraviolet image with a resolution superior to that obtained by ground-based interferometers—the first conventional-telescope image (or "direct-image" in NASA terminology) of the disk of another star.
Because ultraviolet light is absorbed by the Earth's atmosphere, observations at these wavelengths are best performed by space telescopes.
This image, like earlier pictures, contained a bright patch indicating a region in the southwestern quadrant hotter than the stellar surface.
Subsequent ultraviolet spectra taken with the Goddard High Resolution Spectrograph suggested that the hot spot was one of Betelgeuse's poles of rotation. This would give the rotational axis an inclination of about 20° to the direction of Earth, and a position angle from celestial North of about 55°.
2000s studies
In a study published in December 2000, the star's diameter was measured with the Infrared Spatial Interferometer (ISI) at mid-infrared wavelengths producing a limb-darkened estimate of – a figure entirely consistent with Michelson's findings eighty years earlier.
At the time of its publication, the estimated parallax from the Hipparcos mission was , yielding an estimated radius for Betelgeuse of . However, an infrared interferometric study published in 2009 announced that the star had shrunk by 15% since 1993 at an increasing rate without a significant diminution in magnitude.
Subsequent observations suggest that the apparent contraction may be due to shell activity in the star's extended atmosphere.
In addition to the star's diameter, questions have arisen about the complex dynamics of Betelgeuse's extended atmosphere. The mass that makes up galaxies is recycled as stars are formed and destroyed, and red supergiants are major contributors, yet the process by which mass is lost remains a mystery.
With advances in interferometric methodologies, astronomers may be close to resolving this conundrum. Images released by the European Southern Observatory in July 2009, taken by the ground-based Very Large Telescope Interferometer (VLTI), showed a vast plume of gas extending from the star into the surrounding atmosphere.
This mass ejection was equal to the distance between the Sun and Neptune and is one of multiple events occurring in Betelgeuse's surrounding atmosphere. Astronomers have identified at least six shells surrounding Betelgeuse. Solving the mystery of mass loss in the late stages of a star's evolution may reveal those factors that precipitate the explosive deaths of these stellar giants.
2019–2020 fading
A pulsating semiregular variable star, Betelgeuse is subject to multiple cycles of increasing and decreasing brightness due to changes in its size and temperature. The astronomers who first noted the dimming of Betelgeuse, Villanova University astronomers Richard Wasatonic and Edward Guinan, and amateur Thomas Calderwood, theorize that a coincidence of a normal 5.9 year light-cycle minimum and a deeper-than-normal 425 day period are the driving factors.
Other possible causes hypothesized by late 2019 were an eruption of gas or dust or fluctuations in the star's surface brightness.
By August 2020, long-term and extensive studies of Betelgeuse, primarily using ultraviolet observations by the Hubble Space Telescope, had suggested that the unexpected dimming was probably caused by an immense amount of superhot material ejected into space. The material cooled and formed a dust cloud that blocked the starlight coming from about a quarter of Betelgeuse's surface. Hubble captured signs of dense, heated material moving through the star's atmosphere in September, October and November before several telescopes observed the more marked dimming in December and the first few months of 2020.
By January 2020, Betelgeuse had dimmed by a factor of approximately 2.5 from magnitude 0.5 to 1.5 and was reported still fainter in February in The Astronomer's Telegram at a record minimum of +1.614, noting that the star is currently the "least luminous and coolest" in the 25 years of their studies and also calculating a decrease in radius.
Astronomy magazine described it as a "bizarre dimming",
and popular speculation inferred that this might indicate an imminent supernova.
This dropped Betelgeuse from one of the top 10 brightest stars in the sky to outside the top 20, noticeably dimmer than its near neighbor Aldebaran. Mainstream media reports discussed speculation that Betelgeuse might be about to explode as a supernova,
but astronomers note that the supernova is expected to occur within approximately the next 100,000 years and is thus unlikely to be imminent.
By 17 February 2020, Betelgeuse's brightness had remained constant for about 10 days, and the star showed signs of rebrightening.
On 22 February 2020, Betelgeuse may have stopped dimming altogether, all but ending the dimming episode.
On 24 February 2020, no significant change in the infrared over the last 50 years was detected; this seemed unrelated to the recent visual fading and suggested that an impending core collapse may be unlikely.
Also on 24 February 2020, further studies suggested that occluding "large-grain circumstellar dust" may be the most likely explanation for the dimming of the star.
A study that uses observations at submillimetre wavelengths rules out significant contributions from dust absorption. Instead, large starspots appear to be the cause for the dimming.
Followup studies, reported on 31 March 2020 in The Astronomer's Telegram, found a rapid rise in the brightness of Betelgeuse.
Betelgeuse is almost unobservable from the ground between May and August because it is too close to the Sun. Before entering its 2020 conjunction with the Sun, Betelgeuse had reached a brightness of +0.4 . Observations with the STEREO-A spacecraft made in June and July 2020 showed that the star had dimmed by 0.5 since the last ground-based observation in April. This is surprising, because a maximum was expected for August/September 2020, and the next minimum should occur around April 2021. However Betelgeuse's brightness is known to vary irregularly, making predictions difficult. The fading could indicate that another dimming event might occur much earlier than expected.
On 30 August 2020, astronomers reported the detection of a second dust cloud emitted from Betelgeuse, and associated with recent substantial dimming (a secondary minimum on 3 August) in luminosity of the star.
In June 2021, the dust was explained as possibly caused by a cool patch on its photosphere
and in August a second independent group confirmed these results. The dust is thought to have resulted from the cooling of gas ejected from the star. An August 2022
study using the Hubble Space Telescope confirmed previous research and suggested the dust could have been created by a surface mass ejection. It conjectured as well that the dimming could have come from a short-term minimum coinciding with a long-term minimum producing a grand minimum, a 416-day cycle and 2010 day cycle respectively, a mechanism first suggested by astronomer L. Goldberg.
In April 2023, astronomers reported the star reached a peak of 0.0 visual and 0.1 V-band magnitude.
Observation
As a result of its distinctive orange-red color and position within Orion, Betelgeuse is easy to find with the naked eye. It is one of three stars that make up the Winter Triangle asterism, and it marks the center of the Winter Hexagon. It can be seen rising in the east at the beginning of January of each year, just after sunset. Between mid-September and mid-March (best in mid-December), it is visible to virtually every inhabited region of the globe, except in Antarctica at latitudes south of 82°. In May (moderate northern latitudes) or June (southern latitudes), the red supergiant can be seen briefly on the western horizon after sunset, reappearing again a few months later on the eastern horizon before sunrise. In the intermediate period (June–July, centered around mid June), it is invisible to the naked eye (visible only with a telescope in daylight), except around midday low in the north in Antarctic regions between 70° and 80° south latitude (during midday twilight in polar night, when the Sun is below the horizon).
Betelgeuse is a variable star whose visual magnitude ranges between 0.0 and +1.6 . There are periods during which it surpasses Rigel to become the sixth brightest star, and occasionally it will become even brighter than Capella. At its faintest, Betelgeuse can fall behind Deneb and Beta Crucis, themselves both slightly variable, to be the twentieth-brightest star.
Betelgeuse has a B–V color index of 1.85 – a figure which points to its pronounced "redness". The photosphere has an extended atmosphere, which displays strong lines of emission rather than absorption, a phenomenon that occurs when a star is surrounded by a thick gaseous envelope (rather than ionized). This extended gaseous atmosphere has been observed moving toward and away from Betelgeuse, depending on fluctuations in the photosphere. Betelgeuse is the brightest near-infrared source in the sky with a J band magnitude of −2.99; only about 13% of the star's radiant energy is emitted as visible light. If human eyes were sensitive to radiation at all wavelengths, Betelgeuse would appear as the brightest star in the night sky.
Catalogues list up to nine faint visual companions to Betelgeuse. They are at distances of about one to four arc-minutes and all are fainter than 10th magnitude.
Star system
Betelgeuse generally has been considered to be a single isolated star and a runaway star, not currently associated with any cluster or star-forming region, although its birthplace is unclear. However, starting in 1985, three studies have proposed companion stars to Betelgeuse. In 1985, analysis of polarization data from 1968 through 1983 indicated a close companion with a periodic orbit of about 2.1 years, and by using speckle interferometry, the team concluded that the closer of the two companions was located at (≈9 AU) from the main star with a position angle of 273°, an orbit that would potentially place it within the star's chromosphere. The more distant companion was at (≈77 AU) with a position angle of 278°. Other studies have found no evidence for these companions or have actively refuted their existence, but the possibility of a close companion contributing to the overall flux has never been fully ruled out. High-resolution interferometry of Betelgeuse and its vicinity, far beyond the technology of the 1980s and 1990s, has not detected any companions.
More recently, in 2024, two studies found evidence for a companion star.
One study found that a not yet directly-observed, dust-modulating star or white dwarf of at a distance of AU would be the most likely solution for Betelgeuse's 2170-day secondary periodicity, fluctuating radial velocity, moderate radius and low variation in effective temperature, as of 2024. A second study produced by a different group of researchers showed similar findings.
Distance measurements
Parallax is the apparent change of the position of an object, measured in seconds of arc, caused by the change of position of the observer of that object. Parallax is used in astronomy to estimate distances to the nearest stars. As the Earth orbits the Sun, every star is seen to shift by a fraction of an arc second, which measure, combined with the baseline provided by the Earth's orbit gives the distance to that star. Since the first successful parallax measurement by Friedrich Bessel in 1838, astronomers have been puzzled by Betelgeuse's apparent distance. Knowledge of the star's distance improves the accuracy of other stellar parameters, such as luminosity that, when combined with an angular diameter, can be used to calculate the physical radius and effective temperature; luminosity and isotopic abundances can also be used to estimate the stellar age and mass.
When the first interferometric studies were performed on the star's diameter in 1920, the assumed parallax was . This equated to a distance of or roughly , producing not only an inaccurate radius for the star but every other stellar characteristic. Since then, there has been ongoing work to measure the distance of Betelgeuse, with proposed distances as high as or about .
Before the publication of the Hipparcos Catalogue (1997), there were two slightly conflicting parallax measurements for Betelgeuse. The first, in 1991, gave a parallax of , yielding a distance of roughly or . The second was the Hipparcos Input Catalogue (1993) with a trigonometric parallax of , a distance of or . Given this uncertainty, researchers were adopting a wide range of distance estimates, leading to significant variances in the calculation of the star's attributes.
The results from the Hipparcos mission were released in 1997. The measured parallax of Betelgeuse was , which equated to a distance of roughly or , and had a smaller reported error than previous measurements. However, later evaluation of the Hipparcos parallax measurements for variable stars like Betelgeuse found that the uncertainty of these measurements had been underestimated. In 2007, an improved figure of was calculated, hence a much tighter error factor yielding a distance of roughly or .
In 2008, measurements using the Very Large Array (VLA) produced a radio solution of , equaling a distance of or . As the researcher, Harper, points out: "The revised Hipparcos parallax leads to a larger distance () than the original; however, the astrometric solution still requires a significant cosmic noise of 2.4 mas. Given these results it is clear that the Hipparcos data still contain systematic errors of unknown origin." Although the radio data also have systematic errors, the Harper solution combines the datasets in the hope of mitigating such errors. An updated result from further observations with ALMA and e-Merlin gives a parallax of mas and a distance of pc or ly.
In 2020, new observational data from the space-based Solar Mass Ejection Imager aboard the Coriolis satellite and three different modeling techniques produced a refined parallax of mas, a radius of , and a distance of pc or ly, which, if accurate, would mean Betelgeuse is nearly 25% smaller and 25% closer to Earth than previously thought.
The European Space Agency's current Gaia mission is unable to produce good parallax results for stars like Betelgeuse which are brighter than the approximately V=6 saturation limit of the mission's instruments. Because of this limitation, there was no data on Betelgeuse in Gaia Data Release 2, from 2018 or Data Release 3 from 2022.
Variability
Betelgeuse is classified as a semiregular variable star, indicating that some periodicity is noticeable in the brightness changes, but amplitudes may vary, cycles may have different lengths, and there may be standstills or periods of irregularity. It is placed in subgroup SRc; these are pulsating red supergiants with amplitudes around one magnitude and periods from tens to hundreds of days.
Betelgeuse typically shows only small brightness changes near to magnitude +0.5, although at its extremes it can become as bright as magnitude 0.0 or as faint as magnitude +1.6. Betelgeuse is listed in the General Catalogue of Variable Stars with a possible period of 2,335 days. More detailed analyses have shown a main period near 400 days, a short period of 185 days, and a longer secondary period around 2,100 days. The lowest reliably-recorded V-band magnitude of +1.614 was reported in February 2020.
Radial pulsations of red supergiants are well-modelled and show that periods of a few hundred days are typically due to fundamental and first overtone pulsation. Lines in the spectrum of Betelgeuse show doppler shifts indicating radial velocity changes corresponding, very roughly, to the brightness changes. This demonstrates the nature of the pulsations in size, although corresponding temperature and spectral variations are not clearly seen. Variations in the diameter of Betelgeuse have also been measured directly. First overtone pulsations of 185 days have been observed, and the ratio of the fundamental to overtone periods gives valuable information about the internal structure of the star and its age.
The source of the long secondary periods is unknown, but they cannot be explained by radial pulsations. Interferometric observations of Betelgeuse have shown hotspots that are thought to be created by massive convection cells, a significant fraction of the diameter of the star and each emitting 5–10% of the total light of the star. One theory to explain long secondary periods is that they are caused by the evolution of such cells combined with the rotation of the star. Other theories include close binary interactions, chromospheric magnetic activity influencing mass loss, or non-radial pulsations such as g-modes.
In addition to the discrete dominant periods, small-amplitude stochastic variations are seen. It is proposed that this is due to granulation, similar to the same effect on the sun but on a much larger scale.
Diameter
On 13 December 1920, Betelgeuse became the first star outside the Solar System to have the angular size of its photosphere measured. Although interferometry was still in its infancy, the experiment proved a success. The researchers, using a uniform disk model, determined that Betelgeuse had a diameter of , although the stellar disk was likely 17% larger due to the limb darkening, resulting in an estimate for its angular diameter of about 0.055". Since then, other studies have produced angular diameters that range from 0.042 to . Combining these data with historical distance estimates of 180 to yields a projected radius of the stellar disk of anywhere from 1.2 to . Using the Solar System for comparison, the orbit of Mars is about , Ceres in the asteroid belt , Jupiter —so, assuming Betelgeuse occupying the place of the Sun, its photosphere might extend beyond the Jovian orbit, not quite reaching Saturn at .
The precise diameter has been hard to define for several reasons:
Betelgeuse is a pulsating star, so its diameter changes with time;
The star has no definable "edge" as limb darkening causes the optical emissions to vary in color and decrease the farther one extends out from the center;
Betelgeuse is surrounded by a circumstellar envelope composed of matter ejected from the star—matter which absorbs and emits light—making it difficult to define the photosphere of the star;
Measurements can be taken at varying wavelengths within the electromagnetic spectrum and the difference in reported diameters can be as much as 30–35%, yet comparing one finding with another is difficult as the star's apparent size differs depending on the wavelength used. Studies have shown that the measured angular diameter is considerably larger at ultraviolet wavelengths, decreases through the visible to a minimum in the near-infrared, and increase again in the mid-infrared spectrum;
Atmospheric twinkling limits the resolution obtainable from ground-based telescopes since turbulence degrades angular resolution.
The generally reported radii of large cool stars are Rosseland radii, defined as the radius of the photosphere at a specific optical depth of two-thirds. This corresponds to the radius calculated from the effective temperature and bolometric luminosity. The Rosseland radius differs from directly measured radii, with corrections for limb darkening and the observation wavelength. For example, a measured angular diameter of 55.6 mas would correspond to a Rosseland mean diameter of 56.2 mas, while further corrections for the existence of surrounding dust and gas shells would give a diameter of .
To overcome these challenges, researchers have employed various solutions. Astronomical interferometry, first conceived by Hippolyte Fizeau in 1868, was the seminal concept that has enabled major improvements in modern telescopy and led to the creation of the Michelson interferometer in the 1880s, and the first successful measurement of Betelgeuse. Just as human depth perception increases when two eyes instead of one perceive an object, Fizeau proposed the observation of stars through two apertures instead of one to obtain interferences that would furnish information on the star's spatial intensity distribution. The science evolved quickly and multiple-aperture interferometers are now used to capture speckled images, which are synthesized using Fourier analysis to produce a portrait of high resolution. It was this methodology that identified the hotspots on Betelgeuse in the 1990s. Other technological breakthroughs include adaptive optics, space observatories like Hipparcos, Hubble and Spitzer, and the Astronomical Multi-BEam Recombiner (AMBER), which combines the beams of three telescopes simultaneously, allowing researchers to achieve milliarcsecond spatial resolution.
Observations in different regions of the electromagnetic spectrum—the visible, near-infrared (NIR), mid-infrared (MIR), or radio—produce very different angular measurements. In 1996, Betelgeuse was shown to have a uniform disk of . In 2000, a Space Sciences Laboratory team measured a diameter of , ignoring any possible contribution from hotspots, which are less noticeable in the mid-infrared. Also included was a theoretical allowance for limb darkening, yielding a diameter of . The earlier estimate equates to a radius of roughly or , assuming the 2008 Harper distance of , a figure roughly the size of the Jovian orbit of .
In 2004, a team of astronomers working in the near-infrared announced that the more accurate photospheric measurement was . The study also put forth an explanation as to why varying wavelengths from the visible to mid-infrared produce different diameters: the star is seen through a thick, warm extended atmosphere. At short wavelengths (the visible spectrum) the atmosphere scatters light, thus slightly increasing the star's diameter. At near-infrared wavelengths (K and L bands), the scattering is negligible, so the classical photosphere can be directly seen; in the mid-infrared the scattering increases once more, causing the thermal emission of the warm atmosphere to increase the apparent diameter.
Studies with the IOTA and VLTI published in 2009 brought strong support to the idea of dust shells and a molecular shell (MOLsphere) around Betelgeuse, and yielded diameters ranging from 42.57 to with comparatively insignificant margins of error. In 2011, a third estimate in the near-infrared corroborating the 2009 numbers, this time showing a limb-darkened disk diameter of . The near-infrared photospheric diameter of at the Hipparcos distance of equates to about or . A 2014 paper derives an angular diameter of (equivalent to a uniform disc) using H and K band observations made with the VLTI AMBER instrument.
In 2009 it was announced that the radius of Betelgeuse had shrunk from 1993 to 2009 by 15%, with the 2008 angular measurement equal to . Unlike most earlier papers, this study used measurements at one specific wavelength over 15 years. The diminution in Betelgeuse's apparent size equates to a range of values between seen in 1993 to seen in 2008—a contraction of almost in . The observed contraction is generally believed to be a variation in just a portion of the extended atmosphere around Betelgeuse, and observations at other wavelengths have shown an increase in diameter over a similar period.
The latest models of Betelgeuse adopt a photospheric angular diameter of around , with multiple shells out to 50–. Assuming a distance of , this means a stellar diameter of .
Once considered as having the largest angular diameter of any star in the sky after the Sun, Betelgeuse lost that distinction in 1997 when a group of astronomers measured R Doradus with a diameter of , although R Doradus, being much closer to Earth at about , has a linear diameter roughly one-third that of Betelgeuse.
Occultations
Betelgeuse is too far from the ecliptic to be occulted by the major planets, but those by some asteroids (which are more wide-ranging and much more numerous) occur frequently. A partial occultation by the 19th magnitude asteroid occurred on 2 January 2012. It was partial because the angular diameter of the star was larger than that of the asteroid; the brightness of Betelgeuse dropped by only about 0.01 magnitudes.
The 14th magnitude asteroid 319 Leona was predicted to occult on 12 December 2023, 01:12 UTC. Totality was at first uncertain, and the occulation was projected to only last approximately twelve seconds (visible on a narrow path on Earth's surface, the exact width and location of which was initially uncertain due to lack of precise knowledge of the size and path of the asteroid). Projections were later refined as more data were analyzed for a totality ("ring of fire") of approximately five seconds and a 60 km wide path stretching from Tajikistan, Armenia, Turkey, Greece, Italy, Spain, the Atlantic Ocean, Miami, Florida and the Florida Keys to parts of Mexico. (The serendiptous event would also afford detailed observations of 319 Leona itself.) Among other programmes 80 amateur astronomers in Europe alone have been coordinated by astrophysicist Miguel Montargès, et al. of the Paris Observatory for the event.
Physical characteristics
Betelgeuse is a very large, luminous but cool star classified as an M1-2 Ia-ab red supergiant. The letter "M" in this designation means that it is a red star belonging to the M spectral class and therefore has a relatively low photospheric temperature; the "Ia-ab" suffix luminosity class indicates that it is an intermediate-luminosity supergiant, with properties partway between a normal supergiant and a luminous supergiant. Since 1943, the spectrum of Betelgeuse has served as one of the stable anchor points by which other stars are classified.
Uncertainty in the star's surface temperature, diameter, and distance make it difficult to achieve a precise measurement of Betelgeuse's luminosity, but research from 2012 quotes a luminosity of around , assuming a distance of . Studies since 2001 report effective temperatures ranging from 3,250 to 3,690 K. Values outside this range have previously been reported, and much of the variation is believed to be real, due to pulsations in the atmosphere. The star is also a slow rotator and the most recent velocity recorded was —much slower than Antares which has a rotational velocity of . The rotation period depends on Betelgeuse's size and orientation to Earth, but it has been calculated to take to turn on its axis, inclined at an angle of around to Earth.
In 2004, astronomers using computer simulations speculated that even if Betelgeuse is not rotating it might exhibit large-scale magnetic activity in its extended atmosphere, a factor where even moderately strong fields could have a meaningful influence over the star's dust, wind and mass-loss properties. A series of spectropolarimetric observations obtained in 2010 with the Bernard Lyot Telescope at Pic du Midi Observatory revealed the presence of a weak magnetic field at the surface of Betelgeuse, suggesting that the giant convective motions of supergiant stars are able to trigger the onset of a small-scale dynamo effect.
Mass
Betelgeuse has no known orbital companions, so its mass cannot be calculated by that direct method. Modern mass estimates from theoretical modelling have produced values of , with values of – from older studies. It has been calculated that Betelgeuse began its life as a star of , based on a solar luminosity of 90,000–150,000. A novel method of determining the supergiant's mass was proposed in 2011, arguing for a current stellar mass of with an upper limit of 16.6 and lower of , based on observations of the star's intensity profile from narrow H-band interferometry and using a photospheric measurement of roughly or . A probabilistic age prior analysis give a current mass of and an initial mass of .
Betelgeuse's mass can also be estimated based on its position on the colormagnitudediagram (CMD). Betelgeuse's color may have changed from yellow (or possibly orange; i.e. a yellow supergiant) to red in the last few thousand years, based on a 2022 review of historical records. This color change combined with the CMD suggest a mass of and age of 14 Myr, and a distance from 125 to 150 parsecs (~400 to 500 light years).
Motion
The kinematics of Betelgeuse are complex. The age of Class M supergiants with an initial mass of is roughly 10 million years. Starting from its present position and motion, a projection back in time would place Betelgeuse around farther from the galactic plane—an implausible location, as there is no star formation region there. Moreover, Betelgeuse's projected pathway does not appear to intersect with the 25 Ori subassociation or the far younger Orion Nebula Cluster (ONC, also known as Ori OB1d), particularly since Very Long Baseline Array astrometry yields a distance from Betelgeuse to the ONC of between 389 and . Consequently, it is likely that Betelgeuse has not always had its current motion through space but has changed course at one time or another, possibly the result of a nearby stellar explosion. An observation by the Herschel Space Observatory in January 2013 revealed that the star's winds are crashing against the surrounding interstellar medium.
The most likely star-formation scenario for Betelgeuse is that it is a runaway star from the Orion OB1 association. Originally a member of a high-mass multiple system within Ori OB1a, Betelgeuse was probably formed about 10–12 million years ago, but has evolved rapidly due to its high mass. H. Bouy and J. Alves suggested in 2015 that Betelgeuse may instead be a member of the newly discovered Taurion OB association.
Circumstellar dynamics
In the late phase of stellar evolution, massive stars like Betelgeuse exhibit high rates of mass loss, possibly as much as every , resulting in a complex circumstellar environment that is constantly in flux. In a 2009 paper, stellar mass loss was cited as the "key to understanding the evolution of the universe from the earliest cosmological times to the current epoch, and of planet formation and the formation of life itself". However, the physical mechanism is not well understood. When Martin Schwarzschild first proposed his theory of huge convection cells, he argued it was the likely cause of mass loss in evolved supergiants like Betelgeuse. Recent work has corroborated this hypothesis, yet there are still uncertainties about the structure of their convection, the mechanism of their mass loss, the way dust forms in their extended atmosphere, and the conditions which precipitate their dramatic finale as a type II supernova. In 2001, Graham Harper estimated a stellar wind at every , but research since 2009 has provided evidence of episodic mass loss making any total figure for Betelgeuse uncertain. Current observations suggest that a star like Betelgeuse may spend a portion of its lifetime as a red supergiant, but then cross back across the H–R diagram, pass once again through a brief yellow supergiant phase and then explode as a blue supergiant or Wolf–Rayet star.
Astronomers may be close to solving this mystery. They noticed a large plume of gas extending at least six times its stellar radius indicating that Betelgeuse is not shedding matter evenly in all directions. The plume's presence implies that the spherical symmetry of the star's photosphere, often observed in the infrared, is not preserved in its close environment. Asymmetries on the stellar disk had been reported at different wavelengths. However, due to the refined capabilities of the NACO adaptive optics on the VLT, these asymmetries have come into focus. The two mechanisms that could cause such asymmetrical mass loss, were large-scale convection cells or polar mass loss, possibly due to rotation. Probing deeper with ESO's AMBER, gas in the supergiant's extended atmosphere has been observed vigorously moving up and down, creating bubbles as large as the supergiant itself, leading his team to conclude that such stellar upheaval is behind the massive plume ejection observed by Kervella.
Asymmetric shells
In addition to the photosphere, six other components of Betelgeuse's atmosphere have now been identified. They are a molecular environment otherwise known as the MOLsphere, a gaseous envelope, a chromosphere, a dust environment and two outer shells (S1 and S2) composed of carbon monoxide (CO). Some of these elements are known to be asymmetric while others overlap.
At about 0.45 stellar radii (~2–) above the photosphere, there may lie a molecular layer known as the MOLsphere or molecular environment. Studies show it to be composed of water vapor and carbon monoxide with an effective temperature of about . Water vapor had been originally detected in the supergiant's spectrum in the 1960s with the two Stratoscope projects but had been ignored for decades. The MOLsphere may also contain SiO and Al2O3—molecules which could explain the formation of dust particles.
Another cooler region, the asymmetric gaseous envelope, extends for several radii (~10–) from the photosphere. It is enriched in oxygen and especially in nitrogen relative to carbon. These composition anomalies are likely caused by contamination by CNO-processed material from the inside of Betelgeuse.
Radio-telescope images taken in 1998 confirm that Betelgeuse has a highly complex atmosphere, with a temperature of , similar to that recorded on the star's surface but much lower than surrounding gas in the same region. The VLA images also show this lower-temperature gas progressively cools as it extends outward. Although unexpected, it turns out to be the most abundant constituent of Betelgeuse's atmosphere. "This alters our basic understanding of red-supergiant star atmospheres", explained Jeremy Lim, the team's leader. "Instead of the star's atmosphere expanding uniformly due to gas heated to high temperatures near its surface, it now appears that several giant convection cells propel gas from the star's surface into its atmosphere." This is the same region in which Kervella's 2009 finding of a bright plume, possibly containing carbon and nitrogen and extending at least six photospheric radii in the southwest direction of the star, is believed to exist.
The chromosphere was directly imaged by the Faint Object Camera on board the Hubble Space Telescope in ultraviolet wavelengths. The images also revealed a bright area in the southwest quadrant of the disk. The average radius of the chromosphere in 1996 was about 2.2 times the optical disk (~) and was reported to have a temperature no higher than . However, in 2004 observations with the STIS, Hubble's high-precision spectrometer, pointed to the existence of warm chromospheric plasma at least one arcsecond away from the star. At a distance of , the size of the chromosphere could be up to . The observations have conclusively demonstrated that the warm chromospheric plasma spatially overlaps and co-exists with cool gas in Betelgeuse's gaseous envelope as well as with the dust in its circumstellar dust shells.
The first claim of a dust shell surrounding Betelgeuse was put forth in 1977 when it was noted that dust shells around mature stars often emit large amounts of radiation in excess of the photospheric contribution. Using heterodyne interferometry, it was concluded that the red supergiant emits most of its excess radiation from positions beyond 12 stellar radii or roughly the distance of the Kuiper belt at 50 to 60 AU, which depends on the assumed stellar radius. Since then, there have been studies done of this dust envelope at varying wavelengths yielding decidedly different results. Studies from the 1990s have estimated the inner radius of the dust shell anywhere from 0.5 to , or 100 to . These studies point out that the dust environment surrounding Betelgeuse is not static. In 1994, it was reported that Betelgeuse undergoes sporadic decades-long dust production, followed by inactivity. In 1997, significant changes in the dust shell's morphology in one year were noted, suggesting that the shell is asymmetrically illuminated by a stellar radiation field strongly affected by the existence of photospheric hotspots. The 1984 report of a giant asymmetric dust shell () has not been corroborated by recent studies, although another published the same year said that three dust shells were found extending four light-years from one side of the decaying star, suggesting that Betelgeuse sheds its outer layers as it moves.
Although the exact size of the two outer CO shells remains elusive, preliminary estimates suggest that one shell extends from about 1.5 to 4.0 arcseconds and the other expands as far as 7.0 arcseconds. Assuming the Jovian orbit of as the star radius, the inner shell would extend roughly 50 to 150 stellar radii (~300 to ) with the outer one as far as 250 stellar radii (~). The Sun's heliopause is estimated at 100 AU, so the size of this outer shell would be almost fourteen times the size of the Solar System.
Supersonic bow shock
Betelgeuse is travelling through the interstellar medium at a speed of (i.e. ~) creating a bow shock. The shock is not created by the star, but by its powerful stellar wind as it ejects vast amounts of gas into the interstellar medium at a speed of , heating the material surrounding the star, thereby making it visible in infrared light. Because Betelgeuse is so bright, it was only in 1997 that the bow shock was first imaged. The cometary structure is estimated to be at least one parsec wide, assuming a distance of 643 light-years.
Hydrodynamic simulations of the bow shock made in 2012 indicate that it is very young—less than 30,000 years old—suggesting two possibilities: that Betelgeuse moved into a region of the interstellar medium with different properties only recently or that Betelgeuse has undergone a significant transformation producing a changed stellar wind. A 2012 paper, proposed that this phenomenon was caused by Betelgeuse transitioning from a blue supergiant (BSG) to a red supergiant (RSG). There is evidence that in the late evolutionary stage of a star like Betelgeuse, such stars "may undergo rapid transitions from red to blue and vice versa on the Hertzsprung–Russell diagram, with accompanying rapid changes to their stellar winds and bow shocks." Moreover, if future research bears out this hypothesis, Betelgeuse may prove to have traveled close to 200,000 AU as a red supergiant scattering as much as along its trajectory.
Life phases
Betelgeuse is a red supergiant that has evolved from an O-type main-sequence star. After core hydrogen exhaustion, Betelgeuse evolved into a blue supergiant before evolving into its current red supergiant form. Its core will eventually collapse, producing a supernova explosion and leaving behind a compact remnant. The details depend on the exact initial mass and other physical properties of that main sequence star.
Main sequence
The initial mass of Betelgeuse can only be estimated by testing different stellar evolutionary models to match its current observed properties. The unknowns of both the models and the current properties mean that there is considerable uncertainty in Betelgeuse's initial appearance, but its mass is usually estimated to have been in the range of , with modern models finding values of . Its chemical makeup can be reasonably assumed to have been around 70% hydrogen, 28% helium, and 2.4% heavy elements, slightly more metal-rich than the Sun but otherwise similar. The initial rotation rate is more uncertain, but models with slow to moderate initial rotation rates produce the best matches to Betelgeuse's current properties. That main sequence version of Betelgeuse would have been a hot luminous star with a spectral type such as O9V.
A star would take between 11.5 and 15 million years to reach the red supergiant stage, with more rapidly-rotating stars taking the longest. Rapidly-rotating stars take 9.3 million years to reach the red supergiant stage, while stars with slow rotation take only 8.1 million years. These are the best estimates of Betelgeuse's current age, as the time since its zero age main sequence stage is estimated to be 8.0–8.5 million years as a star with no rotation.
After core hydrogen exhaustion
Betelgeuse's time spent as a red supergiant can be estimated by comparing mass loss rates to the observed circumstellar material, as well as the abundances of heavy elements at the surface. Estimates range from 10,000 years to a maximum of 140,000 years. Betelgeuse appears to undergo short periods of heavy mass loss and is a runaway star moving rapidly through space, so comparisons of its current mass loss to the total lost mass are difficult.
The surface of Betelgeuse shows enhancement of nitrogen, relatively low levels of carbon, and a high proportion of 13C relative to 12C, all indicative of a star that has experienced the first dredge-up. However, the first dredge-up occurs soon after a star reaches the red supergiant phase and so this only means that Betelgeuse has been a red supergiant for at least a few thousand years. The best prediction is that Betelgeuse has already spent around 40,000 years as a red supergiant, having left the main sequence perhaps one million years ago.
The current mass can be estimated from evolutionary models from the initial mass and the expected mass lost so far. For Betelgeuse, the total mass lost is predicted to be no more than about , giving a current mass of , considerably higher than estimated by other means such as pulsational properties or limb-darkening models.
All stars more massive than about are expected to end their lives when their cores collapse, typically producing a supernova explosion. Up to about , a type II-P supernova is always produced from the red supergiant stage.
More massive stars can lose mass quickly enough that they evolve towards higher temperatures before their cores can collapse, particularly for rotating stars and models with especially high mass loss rates. These stars can produce type II-L or type IIb supernovae from yellow or blue supergiants, or type I b/c supernovae from Wolf–Rayet stars. Models of rotating stars predict a peculiar type II supernova similar to SN 1987A from a blue supergiant progenitor. On the other hand, non-rotating models predict a type II-P supernova from a red supergiant progenitor.
The time until Betelgeuse explodes depends on the predicted initial conditions and on the estimate of the time already spent as a red supergiant. The total lifetime from the start of the red supergiant phase to core collapse varies from about 300,000 years for a rotating star, 550,000 years for a rotating star, and up to a million years for a non-rotating star. Given the estimated time since Betelgeuse became a red supergiant, estimates of its remaining lifetime range from a "best guess" of under 100,000 years for a non-rotating model to far longer for rotating models or lower-mass stars. Betelgeuse's suspected birthplace in the Orion OB1 association is the location of several previous supernovae. It is believed that runaway stars may be caused by supernovae, and there is strong evidence that OB stars μ Columbae, AE Aurigae, and 53 Arietis all originated from such explosions in Ori OB1 2.2, 2.7, and 4.9 million years ago.
A typical type II-P supernova emits of neutrinos and produces an explosion with a kinetic energy of . As seen from Earth, Betelgeuse as a type II-P supernova would have a peak apparent magnitude somewhere in the range −8 to −12. This would be easily visible in daylight, with a possible brightness up to a significant fraction of the full moon, though likely not exceeding it. This type of supernova would remain at roughly constant brightness for 2–3 months before rapidly dimming. The visible light is produced mainly by the radioactive decay of cobalt, and sustains its brightness due to the increasing transparency of the cooling hydrogen ejected by the supernova.
Media reporting
Due to misunderstandings caused by the 2009 publication of the star's 15% contraction, apparently of its outer atmosphere, Betelgeuse has frequently been the subject of scare stories and rumors suggesting that it will explode within a year, and leading to exaggerated claims about the consequences of such an event. The timing and prevalence of these rumors have been linked to broader misconceptions of astronomy, particularly to doomsday predictions relating to the Mayan calendrical apocalypse. Betelgeuse is not likely to produce a gamma-ray burst and is not close enough for its X-rays, ultraviolet radiation, or ejected material to cause significant effects on Earth.
Following the dimming of Betelgeuse in December 2019, reports appeared in the science and mainstream media that again included speculation that the star might be about to explode as a supernova – even in the face of scientific research that a supernova is not expected for perhaps 100,000 years. Some outlets reported the magnitude as faint as +1.3 as an unusual and interesting phenomenon, like Astronomy magazine, the National Geographic, and the Smithsonian.
Some mainstream media, like The Washington Post, ABC News in Australia, and Popular Science, reported that a supernova was possible but unlikely, whilst other outlets falsely portrayed a supernova as an imminent realistic possibility. CNN, for example, chose the headline "A giant red star is acting weird and scientists think it may be about to explode", while the New York Post declared Betelgeuse as "due for explosive supernova".
Phil Plait, in his Bad Astronomy blog, noting that Betelgeuse's recent behaviour, "[w]hile unusual . . . isn't unprecedented," argued that the star is not likely to explode "for a long, long time." Dennis Overbye of The New York Times agreed that an explosion was not imminent but added that "astronomers are having fun thinking about it."
Following the eventual supernova, a small dense remnant will be left behind, either a neutron star or black hole. Betelgeuse does not seem to have a core massive enough for a black hole, so the remnant will probably be a neutron star of approximately .
Ethnological attributes
Spelling and pronunciation
Betelgeuse has also been spelled Betelgeux and, in German, Beteigeuze (according to Bode). Betelgeux and Betelgeuze were used until the early 20th century, when the spelling Betelgeuse became universal.
Consensus on its pronunciation is weak and is as varied as its spellings:
Oxford English Dictionary and Royal Astronomical Society of Canada
Oxford English Dictionary
Canadian Oxford Dictionary and Webster's Collegiate Dictionary
Martha Evans Martin, The Friendly Stars
The -urz pronunciations are attempts to render the French eu sound; they only work in r-dropping accents.
Etymology
[[File:Al-Sufi's Orion, 1125 Baghdad copy, Doha Museum of Islamic Art Ms 2. 1998. SO.jpg|thumb|An illustration of Orion (horizontally reversed) in al-Sufi's Book of Fixed Stars. Betelgeuze is annotated as Yad al-Jauzā ("Hand of Orion"), one of the proposed etymological origins of its modern name, and also as Mankib al Jauzā ("Shoulder of Orion").]]
Betelgeuse is often mistranslated as "armpit of the central one". In his 1899 work Star-Names and Their Meanings, American amateur naturalist Richard Hinckley Allen stated the derivation was from the , which he claimed degenerated into a number of forms, including Bed Elgueze, Beit Algueze, Bet El-gueze, and Beteigeuze, to the forms Betelgeuse, Betelguese, Betelgueze and Betelgeux. The star was named Beldengeuze in the Alfonsine Tables, and Italian Jesuit priest and astronomer Giovanni Battista Riccioli had called it Bectelgeuze or Bedalgeuze.
Paul Kunitzsch, Professor of Arabic Studies at the University of Munich, refuted Allen's derivation and instead proposed that the full name is a corruption of the Arabic , meaning "the Hand of al-Jauzā'"; i.e., Orion. European mistransliteration into medieval Latin led to the first character y (ﻴ, with two dots underneath) being misread as a b (ﺒ''', with only one dot underneath). During the Renaissance, the star's name was written as ("house of Orion") or , incorrectly thought to mean "armpit of Orion" (a true translation of "armpit" would be , transliterated as ). This led to the modern rendering as Betelgeuse. Other writers have since accepted Kunitzsch's explanation.
The last part of the name, "-elgeuse", comes from the Arabic , a historical Arabic name of the constellation Orion, a feminine name in old Arabian legend, and of uncertain meaning. Because , the root of , means "middle", roughly means "the Central One". The modern Arabic name for Orion is ("the Giant"), although the use of in the star's name has continued. The 17th-century English translator Edmund Chilmead gave it the name Ied Algeuze ("Orion's Hand"), from Christmannus. Other Arabic names recorded include ("the Right Hand"), ("the Arm"), and ("the Shoulder"), all of al-Jauzā, Orion, as .
Other names
Other names for Betelgeuse included the Persian "the Arm", and Coptic "an Armlet". was its Sanskrit name, as part of a Hindu understanding of the constellation as a running antelope or stag. In traditional Chinese astronomy, the name for Betelgeuse is (, the Fourth Star of the constellation of Three Stars) as the Chinese constellation originally referred to the three stars in Orion's Belt. This constellation was ultimately expanded to ten stars, but the earlier name stuck. In Japan, the Taira, or Heike, clan adopted Betelgeuse and its red color as its symbol, calling the star Heike-boshi, (), while the Minamoto, or Genji, clan chose Rigel and its white color. The two powerful families fought a legendary war in Japanese history, the stars seen as facing each other off and only kept apart by the Belt.Hōei Nojiri"Shin seiza jyunrei"p.19
In Tahitian lore, Betelgeuse was one of the pillars propping up the sky, known as Anâ-varu, the pillar to sit by. It was also called Ta'urua-nui-o-Mere "Great festivity in parental yearnings". A Hawaiian term for it was Kaulua-koko ("brilliant red star"). The Lacandon people of Central America knew it as chäk tulix ("red butterfly").
Astronomy writer Robert Burnham Jr. proposed the term padparadaschah, which denotes a rare orange sapphire in India, for the star.
Mythology
With the history of astronomy intimately associated with mythology and astrology before the scientific revolution, the red star, like the planet Mars that derives its name from a Roman war god, has been closely associated with the martial archetype of conquest for millennia, and by extension, the motif of death and rebirth. Other cultures have produced different myths. Stephen R. Wilk has proposed the constellation of Orion could have represented the Greek mythological figure Pelops, who had an artificial shoulder of ivory made for him, with Betelgeuse as the shoulder, its color reminiscent of the reddish yellow sheen of ivory.
Aboriginal people from the Great Victoria Desert of South Australia incorporated Betelgeuse into their oral traditions as the club of Nyeeruna (Orion), which fills with fire-magic and dissipates before returning. This has been interpreted as showing that early Aboriginal observers were aware of the brightness variations of Betelgeuse. The Wardaman people of northern Australia knew the star as Ya-jungin ("Owl Eyes Flicking"), its variable light signifying its intermittent watching of ceremonies led by the Red Kangaroo Leader Rigel. In South African mythology, Betelgeuse was perceived as a lion casting a predatory gaze toward the three zebras represented by Orion's Belt.
In the Americas, Betelgeuse signifies a severed limb of a man-figure (Orion)—the Taulipang of Brazil know the constellation as Zililkawai, a hero whose leg was cut off by his wife, with the variable light of Betelgeuse linked to the severing of the limb. Similarly, the Lakota people of North America see it as a chief whose arm has been severed.
A Sanskrit name for Betelgeuse is ārdrā ("the moist one"), eponymous of the Ardra lunar mansion in Hindu astrology. The Rigvedic God of storms Rudra presided over the star; this association was linked by 19th-century star enthusiast Richard Hinckley Allen to Orion's stormy nature. The constellations in Macedonian folklore represented agricultural items and animals, reflecting their way of life. To them, Betelgeuse was Orach ("the ploughman"), alongside the rest of Orion, which depicted a plough with oxen. The rising of Betelgeuse at around 3 a.m. in late summer and autumn signified the time for village men to go to the fields and plough. To the Inuit, the appearance of Betelgeuse and Bellatrix high in the southern sky after sunset marked the beginning of spring and lengthening days in late February and early March. The two stars were known as Akuttujuuk ("those [two] placed far apart"), referring to the distance between them, mainly to people from North Baffin Island and Melville Peninsula.
The opposed locations of Orion and Scorpius, with their corresponding bright red variable stars Betelgeuse and Antares, were noted by ancient cultures around the world. The setting of Orion and rising of Scorpius signify the death of Orion by the scorpion. In China they signify brothers and rivals Shen and Shang. The Batak of Sumatra marked their New Year with the first new moon after the sinking of Orion's Belt below the horizon, at which point Betelgeuse remained "like the tail of a rooster". The positions of Betelgeuse and Antares at opposite ends of the celestial sky were considered significant, and their constellations were seen as a pair of scorpions. Scorpion days marked as nights that both constellations could be seen.
In popular culture
As one of the brightest and best-known stars, Betelgeuse has featured in many works of fiction. The star's unusual name inspired the title of the 1988 film Beetlejuice, referring to its titular antagonist, and script writer Michael McDowell was impressed by how many people made the connection. In the popular science fiction series The Hitchhiker's Guide to the Galaxy by Douglas Adams, Ford Prefect was from "a small planet somewhere in the vicinity of Betelgeuse."
Two American navy ships were named after the star, both of them World War II vessels, the launched in 1939 and launched in 1944. In 1979, the French supertanker Betelgeuse was moored off Whiddy Island, discharging oil when it exploded, killing 50 people in one of the worst disasters in Ireland's history.
The Dave Matthews Band song "Black and Blue Bird" references the star. The Blur song "Far Out" from their 1994 album Parklife mentions Betelgeuse in its lyrics.
The Philip Larkin poem "The North Ship", found in the collection of the same name, references the star in the section "Above 80° N", which reads:" 'A woman has ten claws,' /
Sang the drunken boatswain; /
Farther than Betelgeuse, /
More brilliant than Orion /
Or the planets Venus and Mars, /
The star flames on the ocean; /
'A woman has ten claws,' /
Sang the drunken boatswain."Humbert Wolfe wrote a poem about Betelgeuse, which was set to music by Gustav Holst.
Table of angular diameter estimates
This table provides a non-exhaustive list of angular measurements conducted since 1920. Also included is a column providing a current range of radii for each study based on Betelgeuse's most recent distance estimate (Harper et al.'') of .
| Physical sciences | Notable stars | null |
10949755 | https://en.wikipedia.org/wiki/Scolopendridae | Scolopendridae | Scolopendridae (or, in older documents, Scolopendridæ) is a family of large centipedes (class Chilopoda).
Description
Nearly all species in this family have four ocelli (simple eyes) on each side of the head and only 21 pairs of legs, but there are exceptions: two scolopendrid species feature more legs (Scolopendropsis bahiensis, with 21 or 23 leg pairs, and S. duplicata, with 39 or 43 leg pairs), and some scolopendrid species are eyeless and blind (e.g., Cormocephalus sagmus, C. pyropygus, and C. delta). Three Asian members of this family, Scolopendra cataracta, Scolopendra paradoxa, and Scolopendra alcyona, are known to show amphibious behaviour. Two other species, Scolopendra hardwickei and Hemiscolopendra marginata, are known to show sexual dimorphism in the composition of their venom.
Genera
Subfamily Otostigminae (Kraepelin, 1903)
Tribe Otostigmini (Kraeplin, 1903)
Alipes Imhoff, 1854
Alluropus Silvestri, 1911
Digitipes Attems, 1930
Edentistoma Tömösváry,1882
Ethmostigmus Pocock, 1898
Otostigmus Porat, 1876
Rhysida Wood, 1862
Tribe Sterropristini (Verhoeff, 1937)
Sterropristes Attems, 1934
Subfamily Scolopendrinae (Leach, 1814)
Arthrorhabdus Pocock, 1891 (= Arthrorhabdinus)
Asanada Meinert, 1885 (= Pseudocryptops)
Asanadopsis Würmli, 1972
Campylostigmus Ribaut, 1923
Notiasemus Koch, 1985
Procrytops Piton, 1940
Psiloscolopendra Kraepelin, 1903
Rhoda Meinert, 1886 (= Pithopus)
Scolopendra Linnaeus, 1758
Scolopendropsis Brandt, 1841
Tonkinodentus Schileyko, 1992
The earliest record of this family is †Cratoraricrus, an extinct genus from the Early Cretaceous of the Crato Formation of Brazil.
| Biology and health sciences | Myriapoda | Animals |
1566231 | https://en.wikipedia.org/wiki/Tar | Tar | Tar is a dark brown or black viscous liquid of hydrocarbons and free carbon, obtained from a wide variety of organic materials through destructive distillation. Tar can be produced from coal, wood, petroleum, or peat.
Mineral products resembling tar can be produced from fossil hydrocarbons, such as petroleum. Coal tar is produced from coal as a byproduct of coke production.
Terminology
"Tar" and "pitch" can be used interchangeably. Asphalt (naturally occurring pitch) may also be called either "mineral tar" or "mineral pitch". There is a tendency to use "tar" for more liquid substances and "pitch" for more solid (viscoelastic) substances. Both "tar" and "pitch" are applied to viscous forms of asphalt, such as the asphalt found in naturally occurring tar pits (e.g., the La Brea Tar Pits in Los Angeles). "Rangoon tar", also known as "Burmese oil" or "Burmese naphtha", is also a form of petroleum. Oil sands, found extensively in Alberta, Canada, and composed of asphalt, are colloquially referred to as "tar sands".
Wood tar
Since prehistoric times wood tar has been used as a water repellent coating for boats, ships, sails, and roofs. In Scandinavia, it was produced as a cash crop. "Peasant Tar" might be named for the district of its production.
Wood tar is still used as an additive in the flavoring of candy, alcohol, and other foods. Wood tar is microbicidal. Producing tar from wood was known in ancient Greece and has probably been used in Scandinavia since the Iron Age. Production and trade in pine-derived tar was a major contributor in the economies of Northern Europe and Colonial America. Its main use was in preserving wooden sailing vessels against rot. For centuries, dating back at least to the 14th century, tar was among Sweden's most important exports. Sweden exported 13,000 barrels of tar in 1615 and 227,000 barrels in the peak year of 1863. The largest user was the Royal Navy of the United Kingdom. Demand for tar declined with the advent of iron and steel ships. Production nearly stopped in the early 20th century. Traditional wooden boats are still sometimes tarred.
The heating (dry distilling) of pine wood causes tar and pitch to drip away from the wood and leave behind charcoal. Birch bark is used to make particularly fine tar, known as "Russian oil", used in Russian leather protection. The by-products of wood tar are turpentine and charcoal. When deciduous tree woods are subjected to destructive distillation, the products are methanol (wood alcohol) and charcoal.
Tar kilns (, , , ) are dry distillation ovens, historically used in Scandinavia for producing tar from wood. They were built close to the forest, from limestone or from more primitive holes in the ground. The bottom is sloped into an outlet hole to allow the tar to pour out. The wood is split into dimensions of a finger, stacked densely, and finally covered tight with earth and moss. If oxygen can enter, the wood might catch fire, and the production would be ruined. On top of this, a fire is stacked and lit. After a few hours, the tar starts to pour out and continues to do so for a few days.
Uses
Tar was used as seal for roofing shingles and tar paper and to seal the hulls of ships and boats. For millennia, wood tar was used to waterproof sails and boats, but today, sails made from inherently waterproof synthetic substances have reduced the demand for tar. Wood tar is still used to seal traditional wooden boats and the roofs of historic, shingle-roofed churches, as well as painting exterior walls of log buildings. Tar is also a general disinfectant. Pine tar oil, or wood tar oil, is used for the surface treatment of wooden shingle roofs, boats, buckets, and tubs and in the medicine, soap, and rubber industries. Pine tar has good penetration on the rough wood. An old wood tar oil recipe for the treatment of wood is one-third each genuine wood tar, balsam turpentine, and boiled or raw linseed oil or Chinese tung oil.
In Finland, wood tar was once considered a panacea reputed to heal "even those cut in twain through their midriff". A Finnish proverb states that "if sauna, vodka and tar won't help, the disease is fatal." Wood tar is used in traditional Finnish medicine because of its microbicidal properties.
Wood tar is also available diluted as tar water, which has numerous uses:
As a flavoring for candies (e.g., Terva Leijona) and alcohol (Terva Viina).
As a spice for food, like meat.
As a scent for saunas. Tar water is mixed into water, which is turned into steam in the sauna.
As an anti-dandruff agent in shampoo.
As a component of cosmetics.
Mixing tar with linseed oil varnish produces tar paint. Tar paint has a translucent brownish hue and can be used to saturate and tone wood and protect it from weather. Tar paint can also be toned with various pigments, producing translucent colors and preserving the wood texture.
Tar was once used for public humiliation, known as tarring and feathering. By pouring hot wood tar onto somebody's bare skin and waiting for it to cool, they would remain stuck in one position. From there, people would attach feathers to the tar, which would remain stuck on the tarred person for the duration of the punishment. That person would then become a public example for the rest of the day.
Pitch was familiar in 9th-century Iraq, derived from petroleum that became accessible from natural fields in the region. It was sometimes used in the construction of baths or in shipbuilding.
Coal tar
Coal tar was formerly one of the products of gasworks. Tar made from coal or petroleum is considered toxic and carcinogenic because of its high benzene content, though coal tar in low concentrations is used as a topical medicine for conditions such as psoriasis. Coal and petroleum tar has a pungent odor.
Coal tar is listed at number 1999 in the United Nations list of dangerous goods.
| Physical sciences | Hydrocarbons | Chemistry |
1566688 | https://en.wikipedia.org/wiki/Rainbow%20lorikeet | Rainbow lorikeet | The rainbow lorikeet (Trichoglossus moluccanus) is a species of parrot found in Australia. It is common along the eastern seaboard, from northern Queensland to South Australia. Its habitat is rainforest, coastal bush and woodland areas. Six taxa traditionally listed as subspecies of the rainbow lorikeet are now treated as separate species (see Taxonomy).
Rainbow lorikeets have been introduced to Perth, Western Australia; Tasmania; Auckland, New Zealand; and Hong Kong.
Taxonomy
The rainbow lorikeet was formally listed in 1788 by the German naturalist Johann Friedrich Gmelin under the binomial name Psittacus moluccanus. Gmelin cited the French polymath Georges-Louis Leclerc, Comte de Buffon who in 1779 had published a description of "La Perruche à Face Bleu" in his Histoire Naturelle des Oiseaux. The species was illustrated as the "Peluche des Moluques" and as the "Perruche d'Amboine". Gmelin was misled and coined the specific epithet moluccanus as he believed the specimens had come from the Moluccas. The type locality was changed to Botany Bay in Australia by Gregory Mathews in 1916. The rainbow lorikeet is now placed in the genus Trichoglossus that was introduced in 1826 by the English naturalist James Francis Stephens.
Two subspecies are recognised:
The rainbow lorikeet has often included the red-collared lorikeet (T. rubritorquis) as a subspecies, but today most major authorities consider it separate. Additionally, a review in 1997 led to the recommendation of splitting off some of the most distinctive taxa from the Lesser Sundas as separate species, these being the scarlet-breasted lorikeet (T. forsteni), the marigold lorikeet (T. capistratus) and the Flores lorikeet (T. weberi). This is increasingly followed by major authorities. In 2019 The rainbow lorikeet in Australia was split into three: rainbow, coconut (T. haematodus) and red-collared lorikeets (T. rubritorquis).
Three syntypes of Trichoglossus novaehollandiae septentrionalis Robinson (Bull. Liverpool Mus., 2, 1900, p.115) are held in the vertebrate zoology collection of National Museums Liverpool at World Museum, with accession numbers NML-VZ 23.7.1900.4, NML-VZ 23.7.1900.4a, and NML-VZ 23.7.1900.4b. The specimens were collected in Cooktown, Queensland, Australia by E. Olive. The specimens came to the Liverpool national collection via purchase from H. C. Robinson.
Description
The rainbow lorikeet is a medium-sized parrot, with the length ranging from including the tail, and the weight varies from . The plumage of the nominate race, as with all subspecies, is very bright and colorful. The head is deep blue with a greenish-yellow nuchal collar, and the rest of the upper parts (wings, back and tail) are green. The chest is orange/yellow. The belly is deep blue, and the thighs and rump are green. In flight a yellow wing-bar contrasts clearly with the red underwing coverts. There is little to visually distinguish between the sexes.
Juveniles have a black beak, which gradually brightens to orange in the adults.
The markings of Trichoglossus moluccanus resemble those of the coconut lorikeet (Trichoglossus haematodus), but with a blue belly and a more orange breast with little or no blue-black barring.
Dimorphism
Unlike the eclectus parrot, rainbow lorikeets do not have any immediately discernible dimorphic traits. Males and females look identical, and surgical sexing by a vet or DNA analysis of a feather is used to determine the sex of an individual.
Behaviour
Rainbow lorikeets often travel together in pairs and occasionally respond to calls to fly as a flock, then disperse again into pairs. Rainbow lorikeet pairs defend their feeding and nesting areas aggressively against other rainbow lorikeets and other bird species. They chase off not only smaller birds, such as the noisy miner and the little wattlebird, but also larger birds such as the Australian magpie.
Diet
Rainbow lorikeets feed mainly on fruit, pollen and nectar, and possess a tongue adapted especially for their particular diet. The end of the tongue is equipped with a papillate appendage adapted to gathering pollen and nectar from flowers. Nectar from eucalyptus is important in Australia, other important nectar sources are Pittosporum, Grevillea, Spathodea campanulata (African tulip-tree), and Metroxylon sagu (sago palm). In Melanesia coconuts are very important food sources, and rainbow lorikeets are important pollinators of these. They also consume the fruits of Ficus, Trema, Muntingia, as well as papaya and mangoes already opened by fruit bats. They also eat crops such as apples, and will raid maize and sorghum. They are also frequent visitors at bird feeders placed in gardens, which supply store-bought nectar, sunflower seeds, and fruits such as apples, grapes and pears.
In many places, including campsites and suburban gardens, wild lorikeets are so used to humans that they can be hand-fed. The Currumbin Wildlife Sanctuary in Queensland, Australia, is noted for its thousands of lorikeets. Around 8am and 4pm each day the birds gather in a huge, noisy flock in the park's main area. Visitors are encouraged to feed them a specially prepared nectar, and the birds will happily settle on people's arms and heads to consume it. Wild rainbow lorikeets can also be hand-fed by visitors at Lone Pine Koala Sanctuary in Brisbane, Queensland, Australia.
Semi-tame lorikeets are common daily visitors in many Sydney backyards, though many people, ignorant of their dietary requirements, feed them bread or bread coated with honey. This is an inadequate source of the nutrients, vitamins and minerals that the rainbow lorikeet requires and can lead to health and feather formation issues in young lorikeets. Packet mixes with a nutritional mix suitable for feeding lorikeets are generally available from vets and pet stores.
Breeding
In southern Australia, breeding usually occurs from late winter to early summer (August to January). Elsewhere in Australia, breeding has been recorded in every month except March, varying from region to region due to changes in food availability and climate. Nesting sites are variable and can include hollows of tall trees such as eucalypts, palm trunks, or overhanging rock. One population in the Admiralty Islands nests in holes in the ground on predator-free islets. Pairs sometimes nest in the same tree with other rainbow lorikeet pairs, or other bird species. The clutch size is between one and three eggs, which are incubated for around 25 days. Incubation duties are carried out by the female alone.
Rainbow lorikeets are mostly monogamous and remain paired for long periods, if not for life.
Status
Overall, the rainbow lorikeet remains widespread and often common. According to the annual Birdlife Australia census, it is the most commonly observed bird in Australia. It is therefore considered to be of least concern by BirdLife International. The status for some localised subspecies is more precarious, with especially T. h. rosenbergii, the Biak lorikeet (which possibly is worthy of treatment as a separate species), being threatened by habitat loss and capture for the parrot trade.
As a pest
Many fruit orchard owners consider them a pest, as they often fly in groups and strip trees containing fresh fruit. In urban areas, the birds create nuisance noise and foul outdoor areas and vehicles with droppings.
The rainbow lorikeet was accidentally released into the southwest of Western Australia near the University of Western Australia in the 1960s and they have since been classified as a pest. They have a major impact there by competing with indigenous bird species, including domination of food sources and competition for increasingly scarce nesting hollows. Bird species such as the purple-crowned lorikeet, the Carnaby's black cockatoo, and the Australian ringneck are adversely affected or displaced.
A feral population was established in New Zealand after a resident of the North Shore, Auckland, illegally released significant numbers of captive-reared birds in the area in the 1990s, which started breeding in the wild. By 1999, a self-sustaining feral population of 150–200 birds had been established in the region, proving that they could survive and adapt to the New Zealand environment. The Department of Conservation, concerned that rainbow lorikeets would outcompete native honeyeaters and by the possible threat to pristine island habitats such as Little Barrier Island, began eradicating the feral population in 2000. The Ministry for Primary Industries Bio-security, in partnership with DOC and regional councils, now manages rainbow lorikeets under the National Interest Pest Response initiative. The aim of the response is to prevent rainbow lorikeets from becoming established in the wild. Late in 2010, five of these birds were discovered living in the Mount Maunganui area. They were fed for a few days before being trapped by a Ministry of Agriculture & Fisheries contractor.
Diseases
Lorikeet paralysis syndrome
A syndrome of uncertain etiology affects rainbow lorikeets every year. Every year in southeast Queensland and northeast New South Wales thousands become paralysed, most significantly, unable to fly or eat. Because this problem is highly seasonal—occurring only October–June and most intensively December–February—it is likely this is a form of plant poisoning. This pattern suggests it is due to the fruits of an unknown plant, which only blooms from the spring to autumn, and most intensively in the summer.
| Biology and health sciences | Psittaciformes | Animals |
1567358 | https://en.wikipedia.org/wiki/Gustnado | Gustnado | A gustnado is a brief, shallow surface-based vortex which forms within the downburst emanating from a thunderstorm. The name is a portmanteau by elision of "gust front tornado", as gustnadoes form due to non-tornadic straight-line wind features in the downdraft (outflow), specifically within the gust front of strong thunderstorms. Gustnadoes tend to be noticed when the vortices loft sufficient debris or form condensation cloud to be visible although it is the wind that makes the gustnado, similarly to tornadoes. As these eddies very rarely connect from the surface to the cloud base, they are very rarely considered as tornadoes. The gustnado has little in common with tornadoes structurally or dynamically in regard to vertical development, intensity, longevity, or formative process—as classic tornadoes are associated with mesocyclones within the inflow (updraft) of the storm, not the outflow.
The average gustnado lasts a few seconds to a few minutes, although there can be several generations and simultaneous swarms. Most have the winds equivalent to an F0 or F1 tornado (up to ), and are commonly mistaken for tornadoes. However, unlike tornadoes, the rotating column of air in a gustnado usually does not extend all the way to the base of the thundercloud. Gustnadoes actually have more in common with (minor) whirlwinds. They are not considered true tornadoes (unless they connect the surface to the ambient cloud base in which case they'd become a landspout) by most meteorologists and are not included in tornado statistics in most areas. Sometimes referred to as spin-up tornadoes, that term more correctly describes the rare tornadic gustnado that connects the surface to the ambient clouded base, or more commonly to the relatively brief but true tornadoes that are associated with a mesovortex.
The most common setting for a gustnado is along the gust front of a severe thunderstorm (by many definitions, containing wind speeds of at least ), along which horizontal shear of the wind may be large. A particularly common location is along the rear-flank gust front of supercell storms. Gustnadoes probably form owing to shear instability associated with the strong horizontal shear; a relative maximum in vertical vorticity must exist in order for shear instability to be present. The bigger question is probably what the dynamical origin(s) of the vertical vorticity is (are), such as the tilting of horizontal vorticity into the vertical or vertical vorticity in the ambient environment that preexists the storm. Along the rear-flank gust front of supercell storms, vertical vorticity very likely has its origins in the upward tilting of vorticity that can occur within descending air in the presence of baroclinity.
While injuries or deaths are rare from gustnadoes, strong ones can cause damage and they are hazardous to drivers. There is some speculation that a gustnado might have been responsible for the collapse of a stage at the Indiana State Fair on August 13, 2011 which killed 7 people and injured 58.
| Physical sciences | Storms | Earth science |
1569600 | https://en.wikipedia.org/wiki/Thermal%20expansion | Thermal expansion | Thermal expansion is the tendency of matter to increase in length, area, or volume, changing its size and density, in response to an increase in temperature (usually excluding phase transitions).
Substances usually contract with decreasing temperature (thermal contraction), with rare exceptions within limited temperature ranges (negative thermal expansion).
Temperature is a monotonic function of the average molecular kinetic energy of a substance. As energy in particles increases, they start moving faster and faster, weakening the intermolecular forces between them and therefore expanding the substance.
When a substance is heated, molecules begin to vibrate and move more, usually creating more distance between themselves.
The relative expansion (also called strain) divided by the change in temperature is called the material's coefficient of linear thermal expansion and generally varies with temperature.
Prediction
If an equation of state is available, it can be used to predict the values of the thermal expansion at all the required temperatures and pressures, along with many other state functions.
Contraction effects (negative expansion)
A number of materials contract on heating within certain temperature ranges; this is usually called negative thermal expansion, rather than "thermal contraction". For example, the coefficient of thermal expansion of water drops to zero as it is cooled to and then becomes negative below this temperature; this means that water has a maximum density at this temperature, and this leads to bodies of water maintaining this temperature at their lower depths during extended periods of sub-zero weather.
Other materials are also known to exhibit negative thermal expansion. Fairly pure silicon has a negative coefficient of thermal expansion for temperatures between about . ALLVAR Alloy 30, a titanium alloy, exhibits anisotropic negative thermal expansion across a wide range of temperatures.
Factors
Unlike gases or liquids, solid materials tend to keep their shape when undergoing thermal expansion.
Thermal expansion generally decreases with increasing bond energy, which also has an effect on the melting point of solids, so high melting point materials are more likely to have lower thermal expansion. In general, liquids expand slightly more than solids. The thermal expansion of glasses is slightly higher compared to that of crystals. At the glass transition temperature, rearrangements that occur in an amorphous material lead to characteristic discontinuities of coefficient of thermal expansion and specific heat. These discontinuities allow detection of the glass transition temperature where a supercooled liquid transforms to a glass.
Absorption or desorption of water (or other solvents) can change the size of many common materials; many organic materials change size much more due to this effect than due to thermal expansion. Common plastics exposed to water can, in the long term, expand by many percent.
Effect on density
Thermal expansion changes the space between particles of a substance, which changes the volume of the substance while negligibly changing its mass (the negligible amount comes from mass–energy equivalence), thus changing its density, which has an effect on any buoyant forces acting on it. This plays a crucial role in convection of unevenly heated fluid masses, notably making thermal expansion partly responsible for wind and ocean currents.
Coefficients
The coefficient of thermal expansion describes how the size of an object changes with a change in temperature. Specifically, it measures the fractional change in size per degree change in temperature at a constant pressure, such that lower coefficients describe lower propensity for change in size. Several types of coefficients have been developed: volumetric, area, and linear. The choice of coefficient depends on the particular application and which dimensions are considered important. For solids, one might only be concerned with the change along a length, or over some area.
The volumetric thermal expansion coefficient is the most basic thermal expansion coefficient, and the most relevant for fluids. In general, substances expand or contract when their temperature changes, with expansion or contraction occurring in all directions. Substances that expand at the same rate in every direction are called isotropic. For isotropic materials, the area and volumetric thermal expansion coefficient are, respectively, approximately twice and three times larger than the linear thermal expansion coefficient.
In the general case of a gas, liquid, or solid, the volumetric coefficient of thermal expansion is given by
The subscript "p" to the derivative indicates that the pressure is held constant during the expansion, and the subscript V stresses that it is the volumetric (not linear) expansion that enters this general definition. In the case of a gas, the fact that the pressure is held constant is important, because the volume of a gas will vary appreciably with pressure as well as temperature. For a gas of low density this can be seen from the ideal gas law.
For various materials
This section summarizes the coefficients for some common materials.
For isotropic materials the coefficients linear thermal expansion α and volumetric thermal expansion αV are related by .
For liquids usually the coefficient of volumetric expansion is listed and linear expansion is calculated here for comparison.
For common materials like many metals and compounds, the thermal expansion coefficient is inversely proportional to the melting point.
In particular, for metals the relation is:
for halides and oxides
In the table below, the range for α is from 10−7 K−1 for hard solids to 10−3 K−1 for organic liquids. The coefficient α varies with the temperature and some materials have a very high variation; see for example the variation vs. temperature of the volumetric coefficient for a semicrystalline polypropylene (PP) at different pressure, and the variation of the linear coefficient vs. temperature for some steel grades (from bottom to top: ferritic stainless steel, martensitic stainless steel, carbon steel, duplex stainless steel, austenitic steel). The highest linear coefficient in a solid has been reported for a Ti-Nb alloy.
(The formula is usually used for solids.)
In solids
When calculating thermal expansion it is necessary to consider whether the body is free to expand or is constrained. If the body is free to expand, the expansion or strain resulting from an increase in temperature can be simply calculated by using the applicable coefficient of thermal expansion.
If the body is constrained so that it cannot expand, then internal stress will be caused (or changed) by a change in temperature. This stress can be calculated by considering the strain that would occur if the body were free to expand and the stress required to reduce that strain to zero, through the stress/strain relationship characterised by the elastic or Young's modulus. In the special case of solid materials, external ambient pressure does not usually appreciably affect the size of an object and so it is not usually necessary to consider the effect of pressure changes.
Common engineering solids usually have coefficients of thermal expansion that do not vary significantly over the range of temperatures where they are designed to be used, so where extremely high accuracy is not required, practical calculations can be based on a constant, average, value of the coefficient of expansion.
Length
Linear expansion means change in one dimension (length) as opposed to change in volume (volumetric expansion).
To a first approximation, the change in length measurements of an object due to thermal expansion is related to temperature change by a coefficient of linear thermal expansion (CLTE). It is the fractional change in length per degree of temperature change. Assuming negligible effect of pressure, one may write:
where is a particular length measurement and is the rate of change of that linear dimension per unit change in temperature.
The change in the linear dimension can be estimated to be:
This estimation works well as long as the linear-expansion coefficient does not change much over the change in temperature , and the fractional change in length is small . If either of these conditions does not hold, the exact differential equation (using ) must be integrated.
Effects on strain
For solid materials with a significant length, like rods or cables, an estimate of the amount of thermal expansion can be described by the material strain, given by and defined as:
where is the length before the change of temperature and is the length after the change of temperature.
For most solids, thermal expansion is proportional to the change in temperature:
Thus, the change in either the strain or temperature can be estimated by:
where
is the difference of the temperature between the two recorded strains, measured in degrees Fahrenheit, degrees Rankine, degrees Celsius, or kelvin, and is the linear coefficient of thermal expansion in "per degree Fahrenheit", "per degree Rankine", "per degree Celsius", or "per kelvin", denoted by , , , or , respectively. In the field of continuum mechanics, thermal expansion and its effects are treated as eigenstrain and eigenstress.
Area
The area thermal expansion coefficient relates the change in a material's area dimensions to a change in temperature. It is the fractional change in area per degree of temperature change. Ignoring pressure, one may write:
where is some area of interest on the object, and is the rate of change of that area per unit change in temperature.
The change in the area can be estimated as:
This equation works well as long as the area expansion coefficient does not change much over the change in temperature , and the fractional change in area is small . If either of these conditions does not hold, the equation must be integrated.
Volume
For a solid, one can ignore the effects of pressure on the material, and the volumetric (or cubical) thermal expansion coefficient can be written:
where is the volume of the material, and is the rate of change of that volume with temperature.
This means that the volume of a material changes by some fixed fractional amount. For example, a steel block with a volume of 1 cubic meter might expand to 1.002 cubic meters when the temperature is raised by 50 K. This is an expansion of 0.2%. If a block of steel has a volume of 2 cubic meters, then under the same conditions, it would expand to 2.004 cubic meters, again an expansion of 0.2%. The volumetric expansion coefficient would be 0.2% for 50 K, or 0.004% K−1.
If the expansion coefficient is known, the change in volume can be calculated
where is the fractional change in volume (e.g., 0.002) and is the change in temperature (50 °C).
The above example assumes that the expansion coefficient did not change as the temperature changed and the increase in volume is small compared to the original volume. This is not always true, but for small changes in temperature, it is a good approximation. If the volumetric expansion coefficient does change appreciably with temperature, or the increase in volume is significant, then the above equation will have to be integrated:
where is the volumetric expansion coefficient as a function of temperature T, and and are the initial and final temperatures respectively.
Isotropic materials
For isotropic materials the volumetric thermal expansion coefficient is three times the linear coefficient:
This ratio arises because volume is composed of three mutually orthogonal directions. Thus, in an isotropic material, for small differential changes, one-third of the volumetric expansion is in a single axis. As an example, take a cube of steel that has sides of length . The original volume will be and the new volume, after a temperature increase, will be
We can easily ignore the terms as ΔL is a small quantity which on squaring gets much smaller and on cubing gets smaller still.
So
The above approximation holds for small temperature and dimensional changes (that is, when and are small), but it does not hold if trying to go back and forth between volumetric and linear coefficients using larger values of . In this case, the third term (and sometimes even the fourth term) in the expression above must be taken into account.
Similarly, the area thermal expansion coefficient is two times the linear coefficient:
This ratio can be found in a way similar to that in the linear example above, noting that the area of a face on the cube is just . Also, the same considerations must be made when dealing with large values of .
Put more simply, if the length of a cubic solid expands from 1.00 m to 1.01 m, then the area of one of its sides expands from 1.00 m2 to 1.02 m2 and its volume expands from 1.00 m3 to 1.03 m3.
Anisotropic materials
Materials with anisotropic structures, such as crystals (with less than cubic symmetry, for example martensitic phases) and many composites, will generally have different linear expansion coefficients in different directions. As a result, the total volumetric expansion is distributed unequally among the three axes. If the crystal symmetry is monoclinic or triclinic, even the angles between these axes are subject to thermal changes. In such cases it is necessary to treat the coefficient of thermal expansion as a tensor with up to six independent elements. A good way to determine the elements of the tensor is to study the expansion by x-ray powder diffraction. The thermal expansion coefficient tensor for the materials possessing cubic symmetry (for e.g. FCC, BCC) is isotropic.
Temperature dependence
Thermal expansion coefficients of solids usually show little dependence on temperature (except at very low temperatures) whereas liquids can expand at different rates at different temperatures. There are some exceptions: for example, cubic boron nitride exhibits significant variation of its thermal expansion coefficient over a broad range of temperatures. Another example is paraffin which in its solid form has a thermal expansion coefficient that is dependent on temperature.
In gases
Since gases fill the entirety of the container which they occupy, the volumetric thermal expansion coefficient at constant pressure, , is the only one of interest.
For an ideal gas, a formula can be readily obtained by differentiation of the ideal gas law, . This yields
where is the pressure, is the molar volume (, with the total number of moles of gas), is the absolute temperature and is equal to the gas constant.
For an isobaric thermal expansion, , so that and the isobaric thermal expansion coefficient is:
which is a strong function of temperature; doubling the temperature will halve the thermal expansion coefficient.
Absolute zero computation
From 1787 to 1802, it was determined by Jacques Charles (unpublished), John Dalton, and Joseph Louis Gay-Lussac that, at constant pressure, ideal gases expanded or contracted their volume linearly (Charles's law) by about 1/273 parts per degree Celsius of temperature's change up or down, between 0° and 100 °C. This suggested that the volume of a gas cooled at about −273 °C would reach zero.
In October 1848, William Thomson, a 24 year old professor of Natural Philosophy at the University of Glasgow, published the paper On an Absolute Thermometric Scale.
In a footnote Thomson calculated that "infinite cold" (absolute zero) was equivalent to −273 °C (he called the temperature in °C as the "temperature of the air thermometers" of the time). This value of "−273" was considered to be the temperature at which the ideal gas volume reaches zero. By considering a thermal expansion linear with temperature (i.e. a constant coefficient of thermal expansion), the value of absolute zero was linearly extrapolated as the negative reciprocal of 0.366/100 °C – the accepted average coefficient of thermal expansion of an ideal gas in the temperature interval 0–100 °C, giving a remarkable consistency to the currently accepted value of −273.15 °C.
In liquids
The thermal expansion of liquids is usually higher than in solids because the intermolecular forces present in liquids are relatively weak and its constituent molecules are more mobile. Unlike solids, liquids have no definite shape and they take the shape of the container. Consequently, liquids have no definite length and area, so linear and areal expansions of liquids only have significance in that they may be applied to topics such as thermometry and estimates of sea level rising due to global climate change. Sometimes, αL is still calculated from the experimental value of αV.
In general, liquids expand on heating, except cold water; below 4 °C it contracts, leading to a negative thermal expansion coefficient. At higher temperatures it shows more typical behavior, with a positive thermal expansion coefficient.
Apparent and absolute
The expansion of liquids is usually measured in a container. When a liquid expands in a vessel, the vessel expands along with the liquid. Hence the observed increase in volume (as measured by the liquid level) is not the actual increase in its volume. The expansion of the liquid relative to the container is called its apparent expansion, while the actual expansion of the liquid is called real expansion or absolute expansion. The ratio of apparent increase in volume of the liquid per unit rise of temperature to the original volume is called its coefficient of apparent expansion. The absolute expansion can be measured by a variety of techniques, including ultrasonic methods.
Historically, this phenomenon complicated the experimental determination of thermal expansion coefficients of liquids, since a direct measurement of the change in height of a liquid column generated by thermal expansion is a measurement of the apparent expansion of the liquid. Thus the experiment simultaneously measures two coefficients of expansion and measurement of the expansion of a liquid must account for the expansion of the container as well. For example, when a flask with a long narrow stem, containing enough liquid to partially fill the stem itself, is placed in a heat bath, the height of the liquid column in the stem will initially drop, followed immediately by a rise of that height until the whole system of flask, liquid and heat bath has warmed through. The initial drop in the height of the liquid column is not due to an initial contraction of the liquid, but rather to the expansion of the flask as it contacts the heat bath first.
Soon after, the liquid in the flask is heated by the flask itself and begins to expand. Since liquids typically have a greater percent expansion than solids for the same temperature change, the expansion of the liquid in the flask eventually exceeds that of the flask, causing the level of liquid in the flask to rise. For small and equal rises in temperature, the increase in volume (real expansion) of a liquid is equal to the sum of the apparent increase in volume (apparent expansion) of the liquid and the increase in volume of the containing vessel. The absolute expansion of the liquid is the apparent expansion corrected for the expansion of the containing vessel.
Examples and applications
The expansion and contraction of the materials must be considered when designing large structures, when using tape or chain to measure distances for land surveys, when designing molds for casting hot material, and in other engineering applications when large changes in dimension due to temperature are expected.
Thermal expansion is also used in mechanical applications to fit parts over one another, e.g. a bushing can be fitted over a shaft by making its inner diameter slightly smaller than the diameter of the shaft, then heating it until it fits over the shaft, and allowing it to cool after it has been pushed over the shaft, thus achieving a 'shrink fit'. Induction shrink fitting is a common industrial method to pre-heat metal components between 150 °C and 300 °C thereby causing them to expand and allow for the insertion or removal of another component.
There exist some alloys with a very small linear expansion coefficient, used in applications that demand very small changes in physical dimension over a range of temperatures. One of these is Invar 36, with expansion approximately equal to 0.6 K−1. These alloys are useful in aerospace applications where wide temperature swings may occur.
Pullinger's apparatus is used to determine the linear expansion of a metallic rod in the laboratory. The apparatus consists of a metal cylinder closed at both ends (called a steam jacket). It is provided with an inlet and outlet for the steam. The steam for heating the rod is supplied by a boiler which is connected by a rubber tube to the inlet. The center of the cylinder contains a hole to insert a thermometer. The rod under investigation is enclosed in a steam jacket. One of its ends is free, but the other end is pressed against a fixed screw. The position of the rod is determined by a micrometer screw gauge or spherometer.
To determine the coefficient of linear thermal expansion of a metal, a pipe made of that metal is heated by passing steam through it. One end of the pipe is fixed securely and the other rests on a rotating shaft, the motion of which is indicated by a pointer. A suitable thermometer records the pipe's temperature. This enables calculation of the relative change in length per degree temperature change.
The control of thermal expansion in brittle materials is a key concern for a wide range of reasons. For example, both glass and ceramics are brittle and uneven temperature causes uneven expansion which again causes thermal stress and this might lead to fracture. Ceramics need to be joined or work in concert with a wide range of materials and therefore their expansion must be matched to the application. Because glazes need to be firmly attached to the underlying porcelain (or other body type) their thermal expansion must be tuned to 'fit' the body so that crazing or shivering do not occur. Good example of products whose thermal expansion is the key to their success are CorningWare and the spark plug. The thermal expansion of ceramic bodies can be controlled by firing to create crystalline species that will influence the overall expansion of the material in the desired direction. In addition or instead the formulation of the body can employ materials delivering particles of the desired expansion to the matrix. The thermal expansion of glazes is controlled by their chemical composition and the firing schedule to which they were subjected. In most cases there are complex issues involved in controlling body and glaze expansion, so that adjusting for thermal expansion must be done with an eye to other properties that will be affected, and generally trade-offs are necessary.
Thermal expansion can have a noticeable effect on gasoline stored in above-ground storage tanks, which can cause gasoline pumps to dispense gasoline which may be more compressed than gasoline held in underground storage tanks in winter, or less compressed than gasoline held in underground storage tanks in summer.
Heat-induced expansion has to be taken into account in most areas of engineering. A few examples are:
Metal-framed windows need rubber spacers.
Rubber tires need to perform well over a range of temperatures, being passively heated or cooled by road surfaces and weather, and actively heated by mechanical flexing and friction.
Metal hot water heating pipes should not be used in long straight lengths.
Large structures such as railways and bridges need expansion joints in the structures to avoid sun kink.
A gridiron pendulum uses an arrangement of different metals to maintain a more temperature stable pendulum length.
A power line on a hot day is droopy, but on a cold day it is tight. This is because the metals expand under heat.
Expansion joints absorb the thermal expansion in a piping system.
Precision engineering nearly always requires the engineer to pay attention to the thermal expansion of the product. For example, when using a scanning electron microscope small changes in temperature such as 1 degree can cause a sample to change its position relative to the focus point.
Liquid thermometers contain a liquid (usually mercury or alcohol) in a tube, which constrains it to flow in only one direction when its volume expands due to changes in temperature.
A bi-metal mechanical thermometer uses a bimetallic strip and bends due to the differing thermal expansion of the two metals.
| Physical sciences | Thermodynamics | Physics |
1570772 | https://en.wikipedia.org/wiki/Paleosol | Paleosol | In geoscience, paleosol (palaeosol in Great Britain and Australia) is an ancient soil that formed in the past. The definition of the term in geology and paleontology is slightly different from its use in soil science.
In geology and paleontology, a paleosol is a former soil preserved by burial underneath either sediments (alluvium or loess) or volcanic deposits (lava flows or volcanic ash), which in the case of older deposits have lithified into rock. In Quaternary geology, sedimentology, paleoclimatology, and geology in general, it is the typical and accepted practice to use the term "paleosol" to designate such "fossil soils" found buried within sedimentary and volcanic deposits exposed in all continents.
In soil science the definition differs slightly: paleosols are soils formed long ago that have no relationship in their chemical and physical characteristics to the present-day climate or vegetation. Such soils are found within extremely old continental cratons, or in small scattered locations in outliers of other ancient rock domains.
Properties
Because of the changes in the Earth's climate over the last 50 million years, soils formed under tropical rainforest (or even savanna) have become exposed to increasingly arid climates which cause former oxisols, ultisols or even alfisols to dry out in such a manner that a very hard crust is formed. This process has occurred so extensively in most parts of Australia as to restrict soil development—the former soil is effectively the parent material for a new soil, but it is so unweatherable that only a very poorly developed soil can exist in present dry climates, especially when they have become much drier during glacial periods in the Quaternary.
In other parts of Australia and in many parts of Africa, drying out of former soils has not been so severe. This has led to large areas of relict podsols in quite dry climates in the far southern inland of Australia (where temperate rainforest was formerly dominant) and to the formation of torrox soils (a suborder of oxisols) in southern Africa. Here, present climates allow, effectively, the maintenance of the old soils in climates under which they could not have formed from the parent material during the Mesozoic and Paleocene.
Paleosols in this sense are always exceedingly infertile soils, containing available phosphorus levels orders of magnitude lower than in temperate regions with younger soils. Ecological studies have shown that this has forced highly specialised evolution amongst Australian flora to obtain minimal nutrient supplies. The fact that soil formation is not occurring makes ecologically sustainable management even more difficult. However, paleosols often contain the most exceptional biodiversity due to the absence of competition.
Taxonomic classification
The record of paleosols extends into the Precambrian in Earth's history, with rare paleosols older than 2.5 billion years. Geology, biology, and the atmosphere all changed significantly over that time, with dramatic shifts at the Great Oxidation Event (2.42 billion years ago) and during the Paleozoic, when complex animals and land plants proliferated.
Consequently, our modern soil classification system cannot be readily applied to paleosols. For example, a modern alfisol—broadly defined as a forest soil—would not have existed prior to the evolution of trees. More problematically, it is specifically defined by chemical properties that would not be preserved in the rock record. While modern soil orders are often used to describe paleosols in a qualitative sense, a paleosol-specific naming scheme has been proposed, although it is only used sporadically in the literature.
Until a paleosol-specific naming scheme is fully adopted, many paleo-pedologists have stuck to using the taxonomic classification of soils provided by the United States Department of Agriculture (USDA). The USDA soil taxonomy attempts to use the measurable properties and objective features within soils to classify them. The methodology developed a hierarchical structure among the different soil taxa, classifying the soils initially at a general level, then assigning soils to progressively more limited subdivisions.
The USDA soil taxonomy does come with drawbacks, including an emphasis on observable features, new nomenclature, and hierarchical organization. The emphasis on observable features can make the soil taxonomy similar in appearance to a legal document. The hierarchical structure cannot be applied more deeply than the order level regarding paleosols. However, despite these drawbacks, the USDA soil taxonomy is still the most comprehensive and influential soil classification system to date. To distinguish and identify paleosols from one another, certain diagnostic horizons and features need to be taken into account. For instance, all paleosols have an A horizon, but histosols have an O horizon above the A horizon.
Identification
Rye & Holland (1998) laid out five criteria for identifying a paleosol. While this was prompted by the need for more stringent identification of Precambrian paleosols, it is applicable to paleosols of any age. The criteria are:
Formed in situ on bedrock,
soft-sediment deformation at the top of the profile, and
up-profile changes in chemistry, texture, and mineralogy consistent with terrestrial weathering processes.
In the field, physical signs of a paleosol include evidence of horizonation (e.g., color and textural changes), bedrock incorporated into a finer overlying lithology (corestones), and evidence of surface processes (e.g., root traces, organic matter, burrows, redox alteration).
Below is a list of soils and some of their diagnostic features that provides a framework for telling these paleosols, or even modern soils, apart:
Entisol (incipient soil)
Horizons (top-to-bottom): A & C
This soil has a very slight degree of soil formation. Original crystalline, metamorphic, or sedimentary features of the parent material experienced little alteration from soil formation. Most are found on young geomorphic surfaces such as flood plains and on steep slopes where erosion removes material as the soil forms. Signs of early successional vegetation of grasses and other herbs and shrubs. Root traces are diagnostic of this type of paleosol because of the small amount of alteration from their parent material in other respects. However, for Entisols of Ordivician age or older, a peak in magnetic susceptibility is indicative of an Entisol.
Inceptisol (young soil)
Horizons: A, sometimes E, Bw, & C
These soils represent a stage of formation beyond Entisols, but not to the degree of development in other soil orders. Typically can be imagined as having a light-colored surface horizon over a moderately weathered subsurface horizon. Forms in low-rolling parts of landscapes in and around steep mountain fronts. Shrubby woodlands of pole trees that form during recolonization of disturbed ground by forests are particularly characteristic of this paleosol. Open woodlands and wooded grasslands are also characteristic of this paleosol.
Andisol (volcanic ash soil)
Horizons: A, Bw, & C
These are soils of volcanic ash of siliceous nature, consisting of bubbles or shards of volcanic glass with a high internal surface area. This soil weathers rapidly to imogolite and smectite. Thus they are highly fertile, rich in organic matter, and have particularly low bulk density. These properties and the aforementioned weathering products typically alter during burial, sometimes to distinctive minerals like celadonite and clinoptilolite. At least 60% recognizable pyroclastic fragments in thin sections are characteristic of this paleosol. This paleosol forms in and around volcanoes.
Histosol (peaty soil)
Horizons: O, A, sometimes Bg, & C
Organic-rich soils with thick peaty horizons, that form in cool, well drained localities or low-lying, permanently waterlogged areas. The primary formation process is accumulation of peat (organic matter), meaning organic matter is produced faster than it can decompose in the soil. The leaching or formation of gley minerals (pyrite or siderite) overprinting prior soil or sedimentary features is associated with peat accumulation.
Spodosol (sandy forest soil)
Horizons: A, E, sometimes Bh, Bs, & C
A subsurface horizon enriched with iron and aluminum oxides or organic matter is characteristic of Spodosols. Displays opaque cements that form distinctive radially cracked, concretionary rims to abundant quartz grains in thin sections. Spodosols form on hilly bedrock or low, rolling quartz-rich sediments. Found principally in humid climates in which clay and soluble salts are dissolved and washed out of the profile and most common in temperate regions. Characteristic vegetation are conifer forests and other kinds of evergreen woody vegetation that can tolerate low nutrient levels and high soil acidity.
Alfisol (fertile forest soil)
Horizons: A, sometimes E, Bt, sometimes Bk, & C
Base-rich forested soils that have a light-colored surface horizon over a clayey subsurface horizon, rich in exchangeable cations. If paleosols contain nodules of carbonate in a horizon deep within the profile, such base saturation can be assumed. If lacking in carbonate nodules, Alfisols can be distinguished by the abundance of base rich clays or by molecular weathering ratios of alumina/bases of less than 2. These soils are not found at the poles or on high mountain tops.
Ultisol (base-poor forest soil)
Horizons: A, sometimes E, Bt, & C
Base-poor forest soils that are similar to Alfisols at first glance. However, Ultisols are more deeply weathered of mineral nutrients. There should not be any calcareous material anywhere within an Ultisol profile and have molecular weathering ratios of alumina/bases of more than 2. Kaolinite and highly weathered aluminous minerals such as gibbsite are common in the profile. Low-base status is attributed to a long formation time. Form mostly on older parts of landscapes, such as rolling hills of bedrock, high alluvial terraces, and plateau tops. Natural vegetation consists of coniferous or hardwood forests.
Oxisol (tropical deeply-weathered soil)
Horizons: A, Bo, sometimes Bv, & C
Deeply weathered soils with texturally uniform profiles. Dominated by kaolinitic clays or other base-poor oxides such as gibbsite or boehmite. Contains molecular weathering ratios of alumina/bases of 10 or more. These soils have deeply weathered mottled horizons. Characteristic of this type of paleosol is a stable microstructure of sand-sized spherical micropeds of iron-stained lay. Very old, often amounting to tens of millions of years. Found on stable continental locations on gentle slopes of plateaus, terraces, and plains. The natural vegetation for Oxisols is a rainforest.
Vertisol (swelling clay soil)
Horizons: A, Bw, & C
These are uniform, thick, clayey soils that have deep, wide cracks. Cracking can produce a hummock-and-swale topography. Mostly composed of smectitic clays. Most Vertisols are found on intermediate to basaltic materials. Found mainly in flat terrain at the foot of gentle slopes. Climate and vegetation are dry and sparse enough that alkaline reactions can be maintained. Vegetation ranges from grassland to open woodland, with wooded grassland being common.
Mollisol (grassland soil)
Horizons: A, sometimes Bt, Bk, sometimes By, & C
Well-developed, base-rich, surface horizon of intimately mixed clay and organic matter. An abundance of fine root traces and crumb ped structures are characteristic of this paleosol. The surface horizon characteristic of this paleosol is created by fine root systems of grassy vegetation and the burrowing activity of many soil invertebrate species. Mollisols are found in low, rolling, or flat country.
Aridisol (desert soil)
Horizons: A, sometimes Bt, Bk, sometimes By, & C
Forms in arid to semi-arid regions, and that lack of rain allows for the creation of shallow calcareous, gypsiferous, or salty horizons. These cements form large nodules or continuous layers. Light-colored, soft, and often vesicular surface horizon. Subsurface horizons are not cemented with any of the aforementioned cements. Mostly found in low-lying areas because steep slopes in arid regions tend to be eroded back to bedrock. Vegetation is sparse and includes prickly shrubs and cacti.
Gelisol (permafrost soil)
Horizons: A, sometimes By, & C
Soils with ground ice or other permafrost features within one meter of the surface. In paleosols, locations of ice can be preserved as clastic dikes, freeze banding, or other deformations created by ground ice. Tillites and other glacigenic deposits are indicative of Gelisols. These soils form under polar desert, tundra, and taiga vegetation. Includes a surprising array of histic epipedons, desert pavements, salic, and calcic horizons.
Many other factors, such as ped structures, such as the presence of blocky, angular or granular peds and fabric type, like clinobimasepic plasmic fabric, are structures that can help one identify if they are dealing with a paleosol. Some of these structures are very helpful when narrowing down the paleosol that is being identified. However, any paleosol should be verified geochemically before use in proxy-based reconstructions; post-deposition alteration processes, such as potassium metasomatism, can change a paleosol's chemistry without dramatically altering its physical appearance.
Applications
Paleoclimate reconstructions
Palaeosols are frequently used as palaeoclimatological tools for gauging the climate in which they formed. Because rates and styles of weathering are dependent on climatic factors, paleosols can be used to reconstruct variables of past climate. Mean annual precipitation (MAP) and air temperature (MAAT) are two commonly-reconstructed variables which, along with seasonality and in conjunction with other paleoenvironmental tools, can be used to describe past terrestrial climates. A suite of paleoclimatic proxies exist and while they vary in focus, many rely on changes in chemical composition throughout a soil profile that occur during weathering, burial, and post-burial processes.
Their use depends on factors such as post-burial alteration, parent material, and soil order; not every proxy is applicable to every paleosol. Most proxies are applicable to Phanerozoic paleosols (not older), as landscape processes changed dramatically after the rise of land plants. Seasonality (the presence and strength of seasons) requires a more nuanced reconstruction approach. Proposed seasonality proxies primarily rely on a soil wetting/drying process, during which pedogenic carbonate can form; like other proxies, this tool is continually being tested and refined.
Paleoatmosphere reconstructions
Soils form in near-constant contact with the atmosphere, so their chemical composition is affected by the composition of the atmosphere through both direct and indirect pathways. The oxidation of paleosols has been used as an indicator of atmospheric oxygen, which has risen over Earth's history. Paleosols have also been used to reconstruct atmospheric carbon dioxide levels, based on modern studies of soil carbon gas exchange, carbon isotopes in pedogenic carbonate nodules, and mass-balance approaches taking multiple atmospheric gases (typically carbon dioxide, oxygen, and methane) into account. These methods are being actively developed in the field of early Earth research.
Paleobotany
Paleosols are an important archive of information about ancient ecosystems and various components of fossil soils can be used to study past plant life. Paleosols often contain ancient plant materials such as pollen grains and phytoliths, a biomineralized form of silica produced by many plants such as grasses. Both pollen and phytolith fossils from different plant species have characteristic shapes that can be traced back to their parent plants. Over long geological time scales, phytoliths may not necessarily be preserved in paleosols due to ability of the poorly crystalline silica to dissolve.
Another indicator of plant community composition in paleosols is the carbon isotopic signature. The ratio of different carbon isotopes in organic matter in paleosols reflects the proportions of plants using C3 photosynthesis, which grow in cooler and wetter climates, versus plants using C4 photosynthesis, which are better adapted to hotter and drier conditions. Other methods for detecting past plant life in paleosols are based on identifying the remains of leaf waxes, which are slow to break down in soils over time.
Paleoseismology
As records of previous Earth surfaces that can be stacked on one another, paleosol sequences are also useful in the field of paleoseismology.
| Physical sciences | Sedimentology | Earth science |
166681 | https://en.wikipedia.org/wiki/Halibut | Halibut | Halibut is the common name for three species of flatfish in the family of right-eye flounders. In some regions, and less commonly, other species of large flatfish are also referred to as halibut.
The word is derived from haly (holy) and butte (flat fish), for its popularity on Catholic holy days. Halibut are demersal fish and are highly regarded as a food fish as well as a sport fish.
Species
A 2018 cladistic analysis based on genetics and morphology showed that the Greenland halibut diverged from a lineage that gave rise to the Atlantic and Pacific halibuts. The common ancestor of all three diverged from a lineage that gave rise to the genus Verasper, comprising the spotted halibut and barfin flounder.
Genus Hippoglossus
Atlantic halibut, Hippoglossus hippoglossus – lives in the North Atlantic
Pacific halibut, Hippoglossus stenolepis – lives in the North Pacific Ocean
Genus Reinhardtius
Greenland halibut, Reinhardtius hippoglossoides – lives in the cold northern Atlantic, northern Pacific, and Arctic Oceans
Physical characteristics
The Pacific and Atlantic halibut are the world's largest flatfish, with debate over which grows larger. Halibut are dark brown on the top side with a white to off-white underbelly and have very small scales invisible to the naked eye embedded in their skin. Halibut are symmetrical at birth with one eye on each side of the head. Then, about six months later, during larval metamorphosis one eye migrates to the other side of the head. The eyes are permanently set once the skull is fully ossified. At the same time, the stationary-eyed side darkens to match the top side, while the other side remains white. This color scheme disguises halibut from above (blending with the ocean floor) and from below (blending into the light from the sky) and is known as countershading.
The IGFA size record for halibut was apparently broken off the waters of Norway in July 2013 by a , fish. This was awaiting certification as of 2013. In July 2014, a Pacific halibut was caught in Glacier Bay, Alaska; this is, however, discounted from records because the halibut was shot and harpooned before being hauled aboard.
Diet
Halibut feed on almost any fish or animal they can fit into their mouths. Juvenile halibut feed on small crustaceans and other bottom-dwelling organisms. Animals found in their stomachs include sand lance, octopus, crab, salmon, hermit crabs, lamprey, sculpin, cod, pollock, herring, and flounder, as well as other halibut. Halibut live at depths ranging from a few meters to hundreds of meters, and although they spend most of their time near the bottom, halibut may move up in the water column to feed. In most ecosystems, the halibut is near the top of the marine food chain. In the North Pacific, common predators are sea lions, killer whales, salmon sharks and humans.
Sex-determining genes
Halibut species vary in sex determination systems. The Atlantic halibut went down a purely XX/XY route, with the male being heterogametic, around 0.9 to 3.8 million years ago. The sex-determining gene for the Atlantic halibut is likely to be gsdf on chromosome 13. The Pacific halibut went down a ZZ/ZW route, with the female being heterogametic, around 4.5 million years ago. The master sex-determining gene of the Pacific halibut is located on chromosome 9 and it is likely to be bmpr1ba. The gene sox2 is likely to play the same role in the Greenland halibut.
Halibut fishery
The North Pacific commercial halibut fishery dates to the late 19th century and today is one of the region's largest and most lucrative. In Canadian and US waters, long-line fishing predominates, using chunks of octopus ("devilfish") or other bait on circle hooks attached at regular intervals to a weighted line that can extend for several miles across the bottom. The fishing vessel retrieves the line after several hours to a day. The effects of long-line gear on habitats are poorly understood, but could include disturbance of sediments, benthic structures, and other structures.
International management is thought to be necessary, because the species occupies waters of the United States, Canada, Russia, and possibly Japan (where the species is known to the Japanese as ohyo), and matures slowly. Halibut do not reproduce until age eight, when about long, so commercial capture below this length prevents breeding and is against US and Canadian regulations supporting sustainability. Pacific halibut fishing is managed by the International Pacific Halibut Commission.
For most of the modern era, halibut fishery operated as a derby. Regulators declared time slots when fishing was open (typically 24–48 hours at a time) and fishermen raced to catch as many pounds as they could within that interval. This approach accommodated unlimited participation in the fishery while allowing regulators to control the quantity of fish caught annually by controlling the number and timing of openings. The approach led to unsafe fishing, as openings were necessarily set before the weather was known, forcing fishermen to leave port regardless of the weather. The approach limited fresh halibut to the markets to several weeks per year when the gluts would push down the price received by fishermen.
Individual fishing quotas
In 1995, US regulators allocated individual fishing quotas (IFQs) to existing fishery participants based on each vessel's documented historical catch. IFQs grant to holders a specific proportion of each year's total allowable catch (TAC). The fishing season is about eight months. The IFQ system improved both safety and product quality by providing a stable flow of fresh halibut to the marketplace. Critics of the program suggest, since holders can sell their quota and the fish are a public resource, the IFQ system gave a public resource to the private sector. The fisheries were managed through a treaty between the United States and Canada per recommendations of the International Pacific Halibut Commission, formed in 1923.
A significant sport fishery in Alaska and British Columbia has emerged, where halibut are prized game and food fish. Sport fisherman use large rods and reels with line, and often bait with herring, large jigs, or whole salmon heads. Halibut are strong and fight strenuously when exposed to air. Smaller fish will usually be pulled on board with a gaff and may be clubbed or even punched in the head to prevent them from thrashing around on the deck. In both commercial and sport fisheries, standard procedure is to shoot or otherwise subdue very large halibut over before landing them.
Overfishing and population decline
The Atlantic halibut has been a major target of fishing since the 1840s with overfishing causing the depletion of the species in the Georges Bank in 1850, then all the way up to the Canadian Arctic in 1866. In the 1940s the American fishing industry collapsed but the Canadian fishing industry remained until there was a decline in Canadian halibut fishery in the 1970s and 1980s. This allowed the halibut population to briefly rebound before collapsing in the 1990s. Since a low point in the early 2000s, the population has rebounded once again and may be stabilizing, but the species is not nearly as abundant in most locations as it was in the early 1800s.
Atlantic halibut population
Currently, Atlantic halibut is managed as two stocks in Canadian waters, which are the Atlantic Continental Shelf stock and the Gulf of St. Lawrence stock. The Atlantic halibut has two other stocks in the Northwest Atlantic, those being the Gulf of Maine-Georges Bank stock controlled by the United States and one controlled by France near the Saint-Pierre and Miquelon Archipelago. The Georges Bank stock is still considered to be depleted and it is listed as a species of concern in the United States. In the two main populations of Atlantic halibut there are many subpopulations, but many have been lost due to patches of extreme overfishing and the populations remain depleted as a whole from what they were in the 1800s.
Pacific and Greenland halibut populations
The Pacific halibut and Greenland halibut have not had this level of fragmentation, and their population is far larger in the United States' waters, with North Pacific halibut and groundfish fisheries extracting the largest volume of catch out of all United States fishery areas. Sometimes the California halibut is mistaken for a subspecies, but they are not, and are not even a true halibut species. In the North Atlantic, observation of migration indicates that there are only two major populations of Greenland halibut that both stretch vast distances. Those populations being the Northeast one stretching from the Kara Sea to Greenland, and the Northwest one stretching from Newfoundland to Baffin Bay. These stocks had been previously thought to be four different populations, but migration has indicated that they are only two different populations, and that fishing has not fragmented them. New research also indicates that the Greenland halibut originally came from the Pacific Ocean and spread into the Arctic Basin when the Bering Strait opened for a second time around 3 million years ago, and thus the Pacific halibut is its closest living relative.
Evolutionary diversification of fragmented populations
In the Atlantic halibut studies have shown that the Atlantic Continental Shelf stock and the Gulf of St. Lawrence stock have begun to differentiate genetically from each other due to low connectivity between populations, low rates of exchange, and subsequent adaptation to local environments. Some adaptations can show up as changes in life-history trait parameters, which can change on a faster time scale than evolution and cause behavioural segregation. This can occur even in areas with enough genetic mixing to prevent genetic divergence. One small but significant observed adaptation difference in the Atlantic halibut has been that the fish in the warmer Scotian Shelf have a faster growth rate than the halibut in the colder southern Grand Banks. The Pacific halibut population remains largely genetically homologous throughout their range, but there is some variation of life-history traits on a geographic gradient. Despite its large range, the populations of Greenland halibut remain largely homogenous due to a lack of barriers for gene flow between its four major populations. There are small differences between subpopulations due to differing environmental factors, such as salinity and temperature gradients, but not to the degree seen in Atlantic halibut, as gene flow and migration continues throughout many different stocks.
As food
Nutrition
Raw Pacific or Atlantic halibut meat is 80% water and 19% protein, with negligible fat and no carbohydrates (table). In a reference amount, raw halibut contains rich content (20% or more of the Daily Value, DV) of protein, selenium (65% DV), phosphorus (34% DV), vitamin D (32% DV), and several B vitamins: niacin, vitamin B6, and vitamin B12 (42–46% DV).
Cooked halibut meat – presumably through the resulting dehydration – has relatively increased protein content and reduced B vitamin content (per 100 grams), while magnesium, phosphorus, and selenium are rich in content.
Food preparation
Halibut yield large fillets from both sides of the fish, with the small round cheeks providing an additional source of meat. Halibut are often boiled, deep-fried or grilled while fresh. Smoking is more difficult with halibut meat than it is with salmon, due to its ultra-low fat content. Eaten fresh, the meat has a clean taste and requires little seasoning. Halibut is noted for its dense and firm texture.
Halibut have historically been an important food source to Alaska Natives and Canadian First Nations, and continue to be a key element to many coastal subsistence economies. Accommodating the competing interests of commercial, sport, and subsistence users is a challenge.
As of 2008, the Atlantic population was so depleted through overfishing that it might be declared an endangered species. According to Seafood Watch, consumers should avoid Atlantic halibut. Most halibut eaten on the East Coast of the United States is from the Pacific.
In 2012, sport fishermen in Cook Inlet reported increased instances of a condition known as "mushy halibut syndrome". The meat of the affected fish has a "jelly-like" consistency. When cooked it does not flake in the normal manner of halibut but rather falls apart. The meat is still perfectly safe to eat but the appearance and consistency are considered unappetizing. The exact cause of the condition is unknown but may be related to a change in diet.
Other species sometimes called "halibut"
Of the same family (Pleuronectidae) as proper halibut
Kamchatka flounder, Atheresthes evermanni – sometimes called "arrowtooth halibut"
Roundnose flounder, Eopsetta grigorjewi – often called "shotted halibut"
Greenland turbot, Reinhardtius hippoglossoides – often called "Greenland halibut"
Spotted halibut, Verasper variegatus
Family Paralichthyidae
California flounder, Paralichthys californicus – sometimes called "California halibut"
Olive flounder, Paralichthys olivaceus – sometimes called "bastard halibut"
Family Psettodidae
Psettodes erumei – sometimes called "Indian halibut"
Family Carangidae (jack family, not a flatfish)
Black pomfret, Parastromateus niger – sometimes called "Australian halibut"
| Biology and health sciences | Acanthomorpha | null |
166689 | https://en.wikipedia.org/wiki/Interferometry | Interferometry | Interferometry is a technique which uses the interference of superimposed waves to extract information. Interferometry typically uses electromagnetic waves and is an important investigative technique in the fields of astronomy, fiber optics, engineering metrology, optical metrology, oceanography, seismology, spectroscopy (and its applications to chemistry), quantum mechanics, nuclear and particle physics, plasma physics, biomolecular interactions, surface profiling, microfluidics, mechanical stress/strain measurement, velocimetry, optometry, and making holograms.
Interferometers are devices that extract information from interference. They are widely used in science and industry for the measurement of microscopic displacements, refractive index changes and surface irregularities. In the case with most interferometers, light from a single source is split into two beams that travel in different optical paths, which are then combined again to produce interference; two incoherent sources can also be made to interfere under some circumstances. The resulting interference fringes give information about the difference in optical path lengths. In analytical science, interferometers are used to measure lengths and the shape of optical components with nanometer precision; they are the highest-precision length measuring instruments in existence. In Fourier transform spectroscopy they are used to analyze light containing features of absorption or emission associated with a substance or mixture. An astronomical interferometer consists of two or more separate telescopes that combine their signals, offering a resolution equivalent to that of a telescope of diameter equal to the largest separation between its individual elements.
Basic principles
Interferometry makes use of the principle of superposition to combine waves in a way that will cause the result of their combination to have some meaningful property that is diagnostic of the original state of the waves. This works because when two waves with the same frequency combine, the resulting intensity pattern is determined by the phase difference between the two waves—waves that are in phase will undergo constructive interference while waves that are out of phase will undergo destructive interference. Waves which are not completely in phase nor completely out of phase will have an intermediate intensity pattern, which can be used to determine their relative phase difference. Most interferometers use light or some other form of electromagnetic wave.
Typically (see Fig. 1, the well-known Michelson configuration) a single incoming beam of coherent light will be split into two identical beams by a beam splitter (a partially reflecting mirror). Each of these beams travels a different route, called a path, and they are recombined before arriving at a detector. The path difference, the difference in the distance traveled by each beam, creates a phase difference between them. It is this introduced phase difference that creates the interference pattern between the initially identical waves. If a single beam has been split along two paths, then the phase difference is diagnostic of anything that changes the phase along the paths. This could be a physical change in the path length itself or a change in the refractive index along the path.
As seen in Fig. 2a and 2b, the observer has a direct view of mirror M1 seen through the beam splitter, and sees a reflected image 2 of mirror M2. The fringes can be interpreted as the result of interference between light coming from the two virtual images 1 and 2 of the original source S. The characteristics of the interference pattern depend on the nature of the light source and the precise orientation of the mirrors and beam splitter. In Fig. 2a, the optical elements are oriented so that 1 and 2 are in line with the observer, and the resulting interference pattern consists of circles centered on the normal to M1 and M'2. If, as in Fig. 2b, M1 and 2 are tilted with respect to each other, the interference fringes will generally take the shape of conic sections (hyperbolas), but if 1 and 2 overlap, the fringes near the axis will be straight, parallel, and equally spaced. If S is an extended source rather than a point source as illustrated, the fringes of Fig. 2a must be observed with a telescope set at infinity, while the fringes of Fig. 2b will be localized on the mirrors.
Use of white light will result in a pattern of colored fringes (see Fig. 3). The central fringe representing equal path length may be light or dark depending on the number of phase inversions experienced by the two beams as they traverse the optical system. (See Michelson interferometer for a discussion of this.)
History
The law of interference of light was described by Thomas Young in his 1803 Bakerian Lecture to the Royal Society of London. In preparation for the lecture, Young performed a double-aperture experiment that demonstrated interference fringes. His interpretation in terms of the interference of waves was rejected by most scientists at the time because of the dominance of Isaac Newton's corpuscular theory of light proposed a century before.
The French engineer Augustin-Jean Fresnel, unaware of Young's results, began working on a wave theory of light and interference and was introduced to François Arago. Between 1816 and 1818, Fresnel and Arago performed interference experiments at the Paris Observatory. During this time, Arago designed and built the first interferometer, using it to measure the refractive index of moist air relative to dry air, which posed a potential problem for astronomical observations of star positions. The success of Fresnel's wave theory of light was established in his prize-winning memoire of 1819 that predicted and measured diffraction patterns. The Arago interferometer was later employed in 1850 by Leon Foucault to measure the speed of light in air relative to water, and it was used again in 1851 by Hippolyte Fizeau to measure the effect of Fresnel drag on the speed of light in moving water.
Jules Jamin developed the first single-beam interferometer (not requiring a splitting aperture as the Arago interferometer did) in 1856. In 1881, the American physicist Albert A. Michelson, while visiting Hermann von Helmholtz in Berlin, invented the interferometer that is named after him, the Michelson Interferometer, to search for effects of the motion of the Earth on the speed of light. Michelson's null results performed in the basement of the Potsdam Observatory outside of Berlin (the horse traffic in the center of Berlin created too many vibrations), and his later more-accurate null results observed with Edward W. Morley at Case College in Cleveland, Ohio, contributed to the growing crisis of the luminiferous ether. Einstein stated that it was Fizeau's measurement of the speed of light in moving water using the Arago interferometer that inspired his theory of the relativistic addition of velocities.
Categories
Interferometers and interferometric techniques may be categorized by a variety of criteria:
Homodyne versus heterodyne detection
In homodyne detection, the interference occurs between two beams at the same wavelength (or carrier frequency). The phase difference between the two beams results in a change in the intensity of the light on the detector. The resulting intensity of the light after mixing of these two beams is measured, or the pattern of interference fringes is viewed or recorded. Most of the interferometers discussed in this article fall into this category.
The heterodyne technique is used for (1) shifting an input signal into a new frequency range as well as (2) amplifying a weak input signal (assuming use of an active mixer). A weak input signal of frequency f1 is mixed with a strong reference frequency f2 from a local oscillator (LO). The nonlinear combination of the input signals creates two new signals, one at the sum f1 + f2 of the two frequencies, and the other at the difference f1 − f2. These new frequencies are called heterodynes. Typically only one of the new frequencies is desired, and the other signal is filtered out of the output of the mixer. The output signal will have an intensity proportional to the product of the amplitudes of the input signals.
The most important and widely used application of the heterodyne technique is in the superheterodyne receiver (superhet), invented in 1917-18 by U.S. engineer Edwin Howard Armstrong and French engineer Lucien Lévy. In this circuit, the incoming radio frequency signal from the antenna is mixed with a signal from a local oscillator (LO) and converted by the heterodyne technique to a lower fixed frequency signal called the intermediate frequency (IF). This IF is amplified and filtered, before being applied to a detector which extracts the audio signal, which is sent to the loudspeaker.
Optical heterodyne detection is an extension of the heterodyne technique to higher (visible) frequencies. While optical heterodyne interferometry is usually done at a single point it is also possible to perform this widefield.
Double path versus common path
A double-path interferometer is one in which the reference beam and sample beam travel along divergent paths. Examples include the Michelson interferometer, the Twyman–Green interferometer, and the Mach–Zehnder interferometer. After being perturbed by interaction with the sample under test, the sample beam is recombined with the reference beam to create an interference pattern which can then be interpreted.
A common-path interferometer is a class of interferometer in which the reference beam and sample beam travel along the same path. Fig. 4 illustrates the Sagnac interferometer, the fibre optic gyroscope, the point diffraction interferometer, and the lateral shearing interferometer. Other examples of common path interferometer include the Zernike phase-contrast microscope, Fresnel's biprism, the zero-area Sagnac, and the scatterplate interferometer.
Wavefront splitting versus amplitude splitting
Wavefront splitting inferometers
A wavefront splitting interferometer divides a light wavefront emerging from a point or a narrow slit (i.e. spatially coherent light) and, after allowing the two parts of the wavefront to travel through different paths, allows them to recombine. Fig. 5 illustrates Young's interference experiment and Lloyd's mirror. Other examples of wavefront splitting interferometer include the Fresnel biprism, the Billet Bi-Lens, diffraction-grating Michelson interferometer, and the Rayleigh interferometer.
In 1803, Young's interference experiment played a major role in the general acceptance of the wave theory of light. If white light is used in Young's experiment, the result is a white central band of constructive interference corresponding to equal path length from the two slits, surrounded by a symmetrical pattern of colored fringes of diminishing intensity. In addition to continuous electromagnetic radiation, Young's experiment has been performed with individual photons, with electrons, and with buckyball molecules large enough to be seen under an electron microscope.
Lloyd's mirror generates interference fringes by combining direct light from a source (blue lines) and light from the source's reflected image (red lines) from a mirror held at grazing incidence. The result is an asymmetrical pattern of fringes. The band of equal path length, nearest the mirror, is dark rather than bright. In 1834, Humphrey Lloyd interpreted this effect as proof that the phase of a front-surface reflected beam is inverted.
Amplitude-splitting inferometers
An amplitude splitting interferometer uses a partial reflector to divide the amplitude of the incident wave into separate beams which are separated and recombined.
The Fizeau interferometer is shown as it might be set up to test an optical flat. A precisely figured reference flat is placed on top of the flat being tested, separated by narrow spacers. The reference flat is slightly beveled (only a fraction of a degree of beveling is necessary) to prevent the rear surface of the flat from producing interference fringes. Separating the test and reference flats allows the two flats to be tilted with respect to each other. By adjusting the tilt, which adds a controlled phase gradient to the fringe pattern, one can control the spacing and direction of the fringes, so that one may obtain an easily interpreted series of nearly parallel fringes rather than a complex swirl of contour lines. Separating the plates, however, necessitates that the illuminating light be collimated. Fig 6 shows a collimated beam of monochromatic light illuminating the two flats and a beam splitter allowing the fringes to be viewed on-axis.
The Mach–Zehnder interferometer is a more versatile instrument than the Michelson interferometer. Each of the well separated light paths is traversed only once, and the fringes can be adjusted so that they are localized in any desired plane. Typically, the fringes would be adjusted to lie in the same plane as the test object, so that fringes and test object can be photographed together. If it is decided to produce fringes in white light, then, since white light has a limited coherence length, on the order of micrometers, great care must be taken to equalize the optical paths or no fringes will be visible. As illustrated in Fig. 6, a compensating cell would be placed in the path of the reference beam to match the test cell. Note also the precise orientation of the beam splitters. The reflecting surfaces of the beam splitters would be oriented so that the test and reference beams pass through an equal amount of glass. In this orientation, the test and reference beams each experience two front-surface reflections, resulting in the same number of phase inversions. The result is that light traveling an equal optical path length in the test and reference beams produces a white light fringe of constructive interference.
The heart of the Fabry–Pérot interferometer is a pair of partially silvered glass optical flats spaced several millimeters to centimeters apart with the silvered surfaces facing each other. (Alternatively, a Fabry–Pérot etalon uses a transparent plate with two parallel reflecting surfaces.) As with the Fizeau interferometer, the flats are slightly beveled. In a typical system, illumination is provided by a diffuse source set at the focal plane of a collimating lens. A focusing lens produces what would be an inverted image of the source if the paired flats were not present, i.e., in the absence of the paired flats, all light emitted from point A passing through the optical system would be focused at point A'. In Fig. 6, only one ray emitted from point A on the source is traced. As the ray passes through the paired flats, it is multiply reflected to produce multiple transmitted rays which are collected by the focusing lens and brought to point A' on the screen. The complete interference pattern takes the appearance of a set of concentric rings. The sharpness of the rings depends on the reflectivity of the flats. If the reflectivity is high, resulting in a high Q factor (i.e., high finesse), monochromatic light produces a set of narrow bright rings against a dark background. In Fig. 6, the low-finesse image corresponds to a reflectivity of 0.04 (i.e., unsilvered surfaces) versus a reflectivity of 0.95 for the high-finesse image.
Fig. 6 illustrates the Fizeau, Mach–Zehnder, and Fabry–Pérot interferometers. Other examples of amplitude splitting interferometer include the Michelson, Twyman–Green, Laser Unequal Path, and Linnik interferometer.
Michelson-Morley
Michelson and Morley (1887) and other early experimentalists using interferometric techniques in an attempt to measure the properties of the luminiferous aether, used monochromatic light only for initially setting up their equipment, always switching to white light for the actual measurements. The reason is that measurements were recorded visually. Monochromatic light would result in a uniform fringe pattern. Lacking modern means of environmental temperature control, experimentalists struggled with continual fringe drift even though the interferometer might be set up in a basement. Since the fringes would occasionally disappear due to vibrations by passing horse traffic, distant thunderstorms and the like, it would be easy for an observer to "get lost" when the fringes returned to visibility. The advantages of white light, which produced a distinctive colored fringe pattern, far outweighed the difficulties of aligning the apparatus due to its low coherence length. This was an early example of the use of white light to resolve the "2 pi ambiguity".
Applications
Physics and astronomy
In physics, one of the most important experiments of the late 19th century was the famous "failed experiment" of Michelson and Morley which provided evidence for special relativity. Recent repetitions of the Michelson–Morley experiment perform heterodyne measurements of beat frequencies of crossed cryogenic optical resonators. Fig 7 illustrates a resonator experiment performed by Müller et al. in 2003. Two optical resonators constructed from crystalline sapphire, controlling the frequencies of two lasers, were set at right angles within a helium cryostat. A frequency comparator measured the beat frequency of the combined outputs of the two resonators. , the precision by which anisotropy of the speed of light can be excluded in resonator experiments is at the 10−17 level.
Michelson interferometers are used in tunable narrow band optical filters and as the core hardware component of Fourier transform spectrometers.
When used as a tunable narrow band filter, Michelson interferometers exhibit a number of advantages and disadvantages when compared with competing technologies such as Fabry–Pérot interferometers or Lyot filters. Michelson interferometers have the largest field of view for a specified wavelength, and are relatively simple in operation, since tuning is via mechanical rotation of waveplates rather than via high voltage control of piezoelectric crystals or lithium niobate optical modulators as used in a Fabry–Pérot system. Compared with Lyot filters, which use birefringent elements, Michelson interferometers have a relatively low temperature sensitivity. On the negative side, Michelson interferometers have a relatively restricted wavelength range and require use of prefilters which restrict transmittance.
Fig. 8 illustrates the operation of a Fourier transform spectrometer, which is essentially a Michelson interferometer with one mirror movable. (A practical Fourier transform spectrometer would substitute corner cube reflectors for the flat mirrors of the conventional Michelson interferometer, but for simplicity, the illustration does not show this.) An interferogram is generated by making measurements of the signal at many discrete positions of the moving mirror. A Fourier transform converts the interferogram into an actual spectrum.
Fig. 9 shows a doppler image of the solar corona made using a tunable Fabry-Pérot interferometer to recover scans of the solar corona at a number of wavelengths near the FeXIV green line. The picture is a color-coded image of the doppler shift of the line, which may be associated with the coronal plasma velocity towards or away from the satellite camera.
Fabry–Pérot thin-film etalons are used in narrow bandpass filters capable of selecting a single spectral line for imaging; for example, the H-alpha line or the Ca-K line of the Sun or stars. Fig. 10 shows an Extreme ultraviolet Imaging Telescope (EIT) image of the Sun at 195 Ångströms (19.5 nm), corresponding to a spectral line of multiply-ionized iron atoms. EIT used multilayer coated reflective mirrors that were coated with alternate layers of a light "spacer" element (such as silicon), and a heavy "scatterer" element (such as molybdenum). Approximately 100 layers of each type were placed on each mirror, with a thickness of around 10 nm each. The layer thicknesses were tightly controlled so that at the desired wavelength, reflected photons from each layer interfered constructively.
The Laser Interferometer Gravitational-Wave Observatory (LIGO) uses two 4-km Michelson–Fabry–Pérot interferometers for the detection of gravitational waves. In this application, the Fabry–Pérot cavity is used to store photons for almost a millisecond while they bounce up and down between the mirrors. This increases the time a gravitational wave can interact with the light, which results in a better sensitivity at low frequencies. Smaller cavities, usually called mode cleaners, are used for spatial filtering and frequency stabilization of the main laser. The first observation of gravitational waves occurred on September 14, 2015.
The Mach–Zehnder interferometer's relatively large and freely accessible working space, and its flexibility in locating the fringes has made it the interferometer of choice for visualizing flow in wind tunnels, and for flow visualization studies in general. It is frequently used in the fields of aerodynamics, plasma physics and heat transfer to measure pressure, density, and temperature changes in gases.
Mach–Zehnder interferometers are also used to study one of the most counterintuitive predictions of quantum mechanics, the phenomenon known as quantum entanglement.
An astronomical interferometer achieves high-resolution observations using the technique of aperture synthesis, mixing signals from a cluster of comparatively small telescopes rather than a single very expensive monolithic telescope.
Early radio telescope interferometers used a single baseline for measurement. Later astronomical interferometers, such as the Very Large Array illustrated in Fig 11, used arrays of telescopes arranged in a pattern on the ground. A limited number of baselines will result in insufficient coverage. This was alleviated by using the rotation of the Earth to rotate the array relative to the sky. Thus, a single baseline could measure information in multiple orientations by taking repeated measurements, a technique called Earth-rotation synthesis. Baselines thousands of kilometers long were achieved using very long baseline interferometry.
Astronomical optical interferometry has had to overcome a number of technical issues not shared by radio telescope interferometry. The short wavelengths of light necessitate extreme precision and stability of construction. For example, spatial resolution of 1 milliarcsecond requires 0.5 μm stability in a 100 m baseline. Optical interferometric measurements require high sensitivity, low noise detectors that did not become available until the late 1990s. Astronomical "seeing", the turbulence that causes stars to twinkle, introduces rapid, random phase changes in the incoming light, requiring data collection rates to be faster than the rate of turbulence. Despite these technical difficulties, three major facilities are now in operation offering resolutions down to the fractional milliarcsecond range. This linked video shows a movie assembled from aperture synthesis images of the Beta Lyrae system, a binary star system approximately 960 light-years (290 parsecs) away in the constellation Lyra, as observed by the CHARA array with the MIRC instrument. The brighter component is the primary star, or the mass donor. The fainter component is the thick disk surrounding the secondary star, or the mass gainer. The two components are separated by 1 milli-arcsecond. Tidal distortions of the mass donor and the mass gainer are both clearly visible.
The wave character of matter can be exploited to build interferometers. The first examples of matter interferometers were electron interferometers, later followed by neutron interferometers. Around 1990 the first atom interferometers were demonstrated, later followed by interferometers employing molecules.
Electron holography is an imaging technique that photographically records the electron interference pattern of an object, which is then reconstructed to yield a greatly magnified image of the original object. This technique was developed to enable greater resolution in electron microscopy than is possible using conventional imaging techniques. The resolution of conventional electron microscopy is not limited by electron wavelength, but by the large aberrations of electron lenses.
Neutron interferometry has been used to investigate the Aharonov–Bohm effect, to examine the effects of gravity acting on an elementary particle, and to demonstrate a strange behavior of fermions that is at the basis of the Pauli exclusion principle: Unlike macroscopic objects, when fermions are rotated by 360° about any axis, they do not return to their original state, but develop a minus sign in their wave function. In other words, a fermion needs to be rotated 720° before returning to its original state.
Atom interferometry techniques are reaching sufficient precision to allow laboratory-scale tests of general relativity.
Interferometers are used in atmospheric physics for high-precision measurements of trace gases via remote sounding of the atmosphere. There are several examples of interferometers that utilize either absorption or emission features of trace gases. A typical use would be in continual monitoring of the column concentration of trace gases such as ozone and carbon monoxide above the instrument.
Engineering and applied science
Newton (test plate) interferometry is frequently used in the optical industry for testing the quality of surfaces as they are being shaped and figured. Fig. 13 shows photos of reference flats being used to check two test flats at different stages of completion, showing the different patterns of interference fringes. The reference flats are resting with their bottom surfaces in contact with the test flats, and they are illuminated by a monochromatic light source. The light waves reflected from both surfaces interfere, resulting in a pattern of bright and dark bands. The surface in the left photo is nearly flat, indicated by a pattern of straight parallel interference fringes at equal intervals. The surface in the right photo is uneven, resulting in a pattern of curved fringes. Each pair of adjacent fringes represents a difference in surface elevation of half a wavelength of the light used, so differences in elevation can be measured by counting the fringes. The flatness of the surfaces can be measured to millionths of an inch by this method. To determine whether the surface being tested is concave or convex with respect to the reference optical flat, any of several procedures may be adopted. One can observe how the fringes are displaced when one presses gently on the top flat. If one observes the fringes in white light, the sequence of colors becomes familiar with experience and aids in interpretation. Finally one may compare the appearance of the fringes as one moves ones head from a normal to an oblique viewing position. These sorts of maneuvers, while common in the optical shop, are not suitable in a formal testing environment. When the flats are ready for sale, they will typically be mounted in a Fizeau interferometer for formal testing and certification.
Fabry-Pérot etalons are widely used in telecommunications, lasers and spectroscopy to control and measure the wavelengths of light. Dichroic filters are multiple layer thin-film etalons. In telecommunications, wavelength-division multiplexing, the technology that enables the use of multiple wavelengths of light through a single optical fiber, depends on filtering devices that are thin-film etalons. Single-mode lasers employ etalons to suppress all optical cavity modes except the single one of interest.
The Twyman–Green interferometer, invented by Twyman and Green in 1916, is a variant of the Michelson interferometer widely used to test optical components. The basic characteristics distinguishing it from the Michelson configuration are the use of a monochromatic point light source and a collimator. Michelson (1918) criticized the Twyman–Green configuration as being unsuitable for the testing of large optical components, since the light sources available at the time had limited coherence length. Michelson pointed out that constraints on geometry forced by limited coherence length required the use of a reference mirror of equal size to the test mirror, making the Twyman–Green impractical for many purposes. Decades later, the advent of laser light sources answered Michelson's objections. (A Twyman–Green interferometer using a laser light source and unequal path length is known as a Laser Unequal Path Interferometer, or LUPI.) Fig. 14 illustrates a Twyman–Green interferometer set up to test a lens. Light from a monochromatic point source is expanded by a diverging lens (not shown), then is collimated into a parallel beam. A convex spherical mirror is positioned so that its center of curvature coincides with the focus of the lens being tested. The emergent beam is recorded by an imaging system for analysis.
Mach–Zehnder interferometers are being used in integrated optical circuits, in which light interferes between two branches of a waveguide that are externally modulated to vary their relative phase. A slight tilt of one of the beam splitters will result in a path difference and a change in the interference pattern. Mach–Zehnder interferometers are the basis of a wide variety of devices, from RF modulators to sensors to optical switches.
The latest proposed extremely large astronomical telescopes, such as the Thirty Meter Telescope and the Extremely Large Telescope, will be of segmented design. Their primary mirrors will be built from hundreds of hexagonal mirror segments. Polishing and figuring these highly aspheric and non-rotationally symmetric mirror segments presents a major challenge. Traditional means of optical testing compares a surface against a spherical reference with the aid of a null corrector. In recent years, computer-generated holograms (CGHs) have begun to supplement null correctors in test setups for complex aspheric surfaces. Fig. 15 illustrates how this is done. Unlike the figure, actual CGHs have line spacing on the order of 1 to 10 μm. When laser light is passed through the CGH, the zero-order diffracted beam experiences no wavefront modification. The wavefront of the first-order diffracted beam, however, is modified to match the desired shape of the test surface. In the illustrated Fizeau interferometer test setup, the zero-order diffracted beam is directed towards the spherical reference surface, and the first-order diffracted beam is directed towards the test surface in such a way that the two reflected beams combine to form interference fringes. The same test setup can be used for the innermost mirrors as for the outermost, with only the CGH needing to be exchanged.
Ring laser gyroscopes (RLGs) and fibre optic gyroscopes (FOGs) are interferometers used in navigation systems. They operate on the principle of the Sagnac effect. The distinction between RLGs and FOGs is that in a RLG, the entire ring is part of the laser while in a FOG, an external laser injects counter-propagating beams into an optical fiber ring, and rotation of the system then causes a relative phase shift between those beams. In a RLG, the observed phase shift is proportional to the accumulated rotation, while in a FOG, the observed phase shift is proportional to the angular velocity.
In telecommunication networks, heterodyning is used to move frequencies of individual signals to different channels which may share a single physical transmission line. This is called frequency division multiplexing (FDM). For example, a coaxial cable used by a cable television system can carry 500 television channels at the same time because each one is given a different frequency, so they don't interfere with one another. Continuous wave (CW) doppler radar detectors are basically heterodyne detection devices that compare transmitted and reflected beams.
Optical heterodyne detection is used for coherent Doppler lidar measurements capable of detecting very weak light scattered in the atmosphere and monitoring wind speeds with high accuracy. It has application in optical fiber communications, in various high resolution spectroscopic techniques, and the self-heterodyne method can be used to measure the linewidth of a laser.
Optical heterodyne detection is an essential technique used in high-accuracy measurements of the frequencies of optical sources, as well as in the stabilization of their frequencies. Until a relatively few years ago, lengthy frequency chains were needed to connect the microwave frequency of a cesium or other atomic time source to optical frequencies. At each step of the chain, a frequency multiplier would be used to produce a harmonic of the frequency of that step, which would be compared by heterodyne detection with the next step (the output of a microwave source, far infrared laser, infrared laser, or visible laser). Each measurement of a single spectral line required several years of effort in the construction of a custom frequency chain. Currently, optical frequency combs have provided a much simpler method of measuring optical frequencies. If a mode-locked laser is modulated to form a train of pulses, its spectrum is seen to consist of the carrier frequency surrounded by a closely spaced comb of optical sideband frequencies with a spacing equal to the pulse repetition frequency (Fig. 16). The pulse repetition frequency is locked to that of the frequency standard, and the frequencies of the comb elements at the red end of the spectrum are doubled and heterodyned with the frequencies of the comb elements at the blue end of the spectrum, thus allowing the comb to serve as its own reference. In this manner, locking of the frequency comb output to an atomic standard can be performed in a single step. To measure an unknown frequency, the frequency comb output is dispersed into a spectrum. The unknown frequency is overlapped with the appropriate spectral segment of the comb and the frequency of the resultant heterodyne beats is measured.
One of the most common industrial applications of optical interferometry is as a versatile measurement tool for the high precision examination of surface topography. Popular interferometric measurement techniques include Phase Shifting Interferometry (PSI), and Vertical Scanning Interferometry(VSI), also known as scanning white light interferometry (SWLI) or by the ISO term coherence scanning interferometry (CSI), CSI exploits coherence to extend the range of capabilities for interference microscopy. These techniques are widely used in micro-electronic and micro-optic fabrication. PSI uses monochromatic light and provides very precise measurements; however it is only usable for surfaces that are very smooth. CSI often uses white light and high numerical apertures, and rather than looking at the phase of the fringes, as does PSI, looks for best position of maximum fringe contrast or some other feature of the overall fringe pattern. In its simplest form, CSI provides less precise measurements than PSI but can be used on rough surfaces. Some configurations of CSI, variously known as Enhanced VSI (EVSI), high-resolution SWLI or Frequency Domain Analysis (FDA), use coherence effects in combination with interference phase to enhance precision.
Phase Shifting Interferometry addresses several issues associated with the classical analysis of static interferograms. Classically, one measures the positions of the fringe centers. As seen in Fig. 13, fringe deviations from straightness and equal spacing provide a measure of the aberration. Errors in determining the location of the fringe centers provide the inherent limit to precision of the classical analysis, and any intensity variations across the interferogram will also introduce error. There is a trade-off between precision and number of data points: closely spaced fringes provide many data points of low precision, while widely spaced fringes provide a low number of high precision data points. Since fringe center data is all that one uses in the classical analysis, all of the other information that might theoretically be obtained by detailed analysis of the intensity variations in an interferogram is thrown away. Finally, with static interferograms, additional information is needed to determine the polarity of the wavefront: In Fig. 13, one can see that the tested surface on the right deviates from flatness, but one cannot tell from this single image whether this deviation from flatness is concave or convex. Traditionally, this information would be obtained using non-automated means, such as by observing the direction that the fringes move when the reference surface is pushed.
Phase shifting interferometry overcomes these limitations by not relying on finding fringe centers, but rather by collecting intensity data from every point of the CCD image sensor. As seen in Fig. 17, multiple interferograms (at least three) are analyzed with the reference optical surface shifted by a precise fraction of a wavelength between each exposure using a piezoelectric transducer (PZT). Alternatively, precise phase shifts can be introduced by modulating the laser frequency. The captured images are processed by a computer to calculate the optical wavefront errors. The precision and reproducibility of PSI is far greater than possible in static interferogram analysis, with measurement repeatabilities of a hundredth of a wavelength being routine. Phase shifting technology has been adapted to a variety of interferometer types such as Twyman–Green, Mach–Zehnder, laser Fizeau, and even common path configurations such as point diffraction and lateral shearing interferometers. More generally, phase shifting techniques can be adapted to almost any system that uses fringes for measurement, such as holographic and speckle interferometry.
In coherence scanning interferometry, interference is only achieved when the path length delays of the interferometer are matched within the coherence time of the light source. CSI monitors the fringe contrast rather than the phase of the fringes. Fig. 17 illustrates a CSI microscope using a Mirau interferometer in the objective; other forms of interferometer used with white light include the Michelson interferometer (for low magnification objectives, where the reference mirror in a Mirau objective would interrupt too much of the aperture) and the Linnik interferometer (for high magnification objectives with limited working distance). The sample (or alternatively, the objective) is moved vertically over the full height range of the sample, and the position of maximum fringe contrast is found for each pixel. The chief benefit of coherence scanning interferometry is that systems can be designed that do not suffer from the 2 pi ambiguity of coherent interferometry, and as seen in Fig. 18, which scans a 180μm x 140μm x 10μm volume, it is well suited to profiling steps and rough surfaces. The axial resolution of the system is determined in part by the coherence length of the light source. Industrial applications include in-process surface metrology, roughness measurement, 3D surface metrology in hard-to-reach spaces and in hostile environments, profilometry of surfaces with high aspect ratio features (grooves, channels, holes), and film thickness measurement (semi-conductor and optical industries, etc.).
Fig. 19 illustrates a Twyman–Green interferometer set up for white light scanning of a macroscopic object.
Holographic interferometry is a technique which uses holography to monitor small deformations in single wavelength implementations. In multi-wavelength implementations, it is used to perform dimensional metrology of large parts and assemblies and to detect larger surface defects.
Holographic interferometry was discovered by accident as a result of mistakes committed during the making of holograms. Early lasers were relatively weak and photographic plates were insensitive, necessitating long exposures during which vibrations or minute shifts might occur in the optical system. The resultant holograms, which showed the holographic subject covered with fringes, were considered ruined.
Eventually, several independent groups of experimenters in the mid-60s realized that the fringes encoded important information about dimensional changes occurring in the subject, and began intentionally producing holographic double exposures. The main Holographic interferometry article covers the disputes over priority of discovery that occurred during the issuance of the patent for this method.
Double- and multi- exposure holography is one of three methods used to create holographic interferograms. A first exposure records the object in an unstressed state. Subsequent exposures on the same photographic plate are made while the object is subjected to some stress. The composite image depicts the difference between the stressed and unstressed states.
Real-time holography is a second method of creating holographic interferograms. A holograph of the unstressed object is created. This holograph is illuminated with a reference beam to generate a hologram image of the object directly superimposed over the original object itself while the object is being subjected to some stress. The object waves from this hologram image will interfere with new waves coming from the object. This technique allows real time monitoring of shape changes.
The third method, time-average holography, involves creating a holograph while the object is subjected to a periodic stress or vibration. This yields a visual image of the vibration pattern.
Interferometric synthetic aperture radar (InSAR) is a radar technique used in geodesy and remote sensing. Satellite synthetic aperture radar images of a geographic feature are taken on separate days, and changes that have taken place between radar images taken on the separate days are recorded as fringes similar to those obtained in holographic interferometry. The technique can monitor centimeter- to millimeter-scale deformation resulting from earthquakes, volcanoes and landslides, and also has uses in structural engineering, in particular for the monitoring of subsidence and structural stability. Fig 20 shows Kilauea, an active volcano in Hawaii. Data acquired using the space shuttle Endeavour's X-band Synthetic Aperture Radar on April 13, 1994 and October 4, 1994 were used to generate interferometric fringes, which were overlaid on the X-SAR image of Kilauea.
Electronic speckle pattern interferometry (ESPI), also known as TV holography, uses video detection and recording to produce an image of the object upon which is superimposed a fringe pattern which represents the displacement of the object between recordings. (see Fig. 21) The fringes are similar to those obtained in holographic interferometry.
When lasers were first invented, laser speckle was considered to be a severe drawback in using lasers to illuminate objects, particularly in holographic imaging because of the grainy image produced. It was later realized that speckle patterns could carry information about the object's surface deformations. Butters and Leendertz developed the technique of speckle pattern interferometry in 1970, and since then, speckle has been exploited in a variety of other applications. A photograph is made of the speckle pattern before deformation, and a second photograph is made of the speckle pattern after deformation. Digital subtraction of the two images results in a correlation fringe pattern, where the fringes represent lines of equal deformation. Short laser pulses in the nanosecond range can be used to capture very fast transient events. A phase problem exists: In the absence of other information, one cannot tell the difference between contour lines indicating a peak versus contour lines indicating a trough. To resolve the issue of phase ambiguity, ESPI may be combined with phase shifting methods.
A method of establishing precise geodetic baselines, invented by Yrjö Väisälä, exploited the low coherence length of white light. Initially, white light was split in two, with the reference beam "folded", bouncing back-and-forth six times between a mirror pair spaced precisely 1 m apart. Only if the test path was precisely 6 times the reference path would fringes be seen. Repeated applications of this procedure allowed precise measurement of distances up to 864 meters. Baselines thus established were used to calibrate geodetic distance measurement equipment, leading to a metrologically traceable scale for geodetic networks measured by these instruments. (This method has been superseded by GPS.)
Other uses of interferometers have been to study dispersion of materials, measurement of complex indices of refraction, and thermal properties. They are also used for three-dimensional motion mapping including mapping vibrational patterns of structures.
Biology and medicine
Optical interferometry, applied to biology and medicine, provides sensitive metrology capabilities for the measurement of biomolecules, subcellular components, cells and tissues. Many forms of label-free biosensors rely on interferometry because the direct interaction of electromagnetic fields with local molecular polarizability eliminates the need for fluorescent tags or nanoparticle markers. At a larger scale, cellular interferometry shares aspects with phase-contrast microscopy, but comprises a much larger class of phase-sensitive optical configurations that rely on optical interference among cellular constituents through refraction and diffraction. At the tissue scale, partially-coherent forward-scattered light propagation through the micro aberrations and heterogeneity of tissue structure provides opportunities to use phase-sensitive gating (optical coherence tomography) as well as phase-sensitive fluctuation spectroscopy to image subtle structural and dynamical properties.
Optical coherence tomography (OCT) is a medical imaging technique using low-coherence interferometry to provide tomographic visualization of internal tissue microstructures. As seen in Fig. 22, the core of a typical OCT system is a Michelson interferometer. One interferometer arm is focused onto the tissue sample and scans the sample in an X-Y longitudinal raster pattern. The other interferometer arm is bounced off a reference mirror. Reflected light from the tissue sample is combined with reflected light from the reference. Because of the low coherence of the light source, interferometric signal is observed only over a limited depth of sample. X-Y scanning therefore records one thin optical slice of the sample at a time. By performing multiple scans, moving the reference mirror between each scan, an entire three-dimensional image of the tissue can be reconstructed. Recent advances have striven to combine the nanometer phase retrieval of coherent interferometry with the ranging capability of low-coherence interferometry.
Phase contrast and differential interference contrast (DIC) microscopy are important tools in biology and medicine. Most animal cells and single-celled organisms have very little color, and their intracellular organelles are almost totally invisible under simple bright field illumination. These structures can be made visible by staining the specimens, but staining procedures are time-consuming and kill the cells. As seen in Figs. 24 and 25, phase contrast and DIC microscopes allow unstained, living cells to be studied. DIC also has non-biological applications, for example in the analysis of planar silicon semiconductor processing.
Angle-resolved low-coherence interferometry (a/LCI) uses scattered light to measure the sizes of subcellular objects, including cell nuclei. This allows interferometry depth measurements to be combined with density measurements. Various correlations have been found between the state of tissue health and the measurements of subcellular objects. For example, it has been found that as tissue changes from normal to cancerous, the average cell nuclei size increases.
Phase-contrast X-ray imaging (Fig. 26) refers to a variety of techniques that use phase information of a coherent x-ray beam to image soft tissues. (For an elementary discussion, see Phase-contrast x-ray imaging (introduction). For a more in-depth review, see Phase-contrast X-ray imaging.) It has become an important method for visualizing cellular and histological structures in a wide range of biological and medical studies. There are several technologies being used for x-ray phase-contrast imaging, all utilizing different principles to convert phase variations in the x-rays emerging from an object into intensity variations. These include propagation-based phase contrast, Talbot interferometry, Moiré-based far-field interferometry, refraction-enhanced imaging, and x-ray interferometry. These methods provide higher contrast compared to normal absorption-contrast x-ray imaging, making it possible to see smaller details. A disadvantage is that these methods require more sophisticated equipment, such as synchrotron or microfocus x-ray sources, x-ray optics, or high resolution x-ray detectors.
| Technology | Optical instruments | null |
166734 | https://en.wikipedia.org/wiki/Altocumulus%20cloud | Altocumulus cloud | Altocumulus () is a middle-altitude cloud genus that belongs mainly to the physical category, characterized by globular masses or rolls in layers or patchesthe individual elements being larger and darker than those of cirrocumulus and smaller than those of stratocumulus. However, if the layers become tufted in appearance due to increased airmass instability, then the altocumulus clouds become more purely cumuliform in structure. Like other cumuliform and stratocumuliform clouds, altocumulus signifies convection. A sheet of partially conjoined altocumulus perlucidus is sometimes found preceding a weakening warm front, where the altostratus is starting to fragment, resulting in patches of altocumulus perlucidus between the areas of altostratus. Altocumulus is also commonly found between the warm and cold fronts in a depression, although this is often hidden by lower clouds.
Towering altocumulus, known as altocumulus castellanus, frequently signals the development of thunderstorms later in the day, as it shows instability and convection in the middle levels of the troposphere, the area where towering cumulus clouds can turn into cumulonimbus. It is therefore one of three warning clouds often recorded by the aviation industry, the other two being towering cumulus and cumulonimbus. Altocumulus generally forms at about above ground level, a similar level to altostratus formations, and satellite photography reveals that the two types of cloud can create formations that can stretch for thousands of square miles. Extensive altocumulus formations, particularly if they take the form of undulatus, are often referred to as altocumulus mackerel sky.
Subtypes
Species: Altocumulus has four species. The stratiformis species (Ac str) is composed of sheets or relatively flat patches of stratocumuliform cloud. The synoptic coding is determined by the predominant variety or occasionally by the genitus mother cloud. Altocumulus lenticularis (Ac len lenticular cloud) is a lens-shaped middle cloud which can resemble flying saucers and may occasionally be mistaken for "unidentified flying objects". This is formed by uplift usually associated with mountains. but usually with at least some grey shading. It is coded CM4 on the SYNOP weather observation. Grey shading is also seen with altocumulus castellanus (Ac cas), a turreted middle cloud that can achieve significant vertical development and signals increasing air mass instability. It is nevertheless usually classified as middle rather than vertical and is coded CM8. The floccus species (Ac flo) is a tufted middle cloud which is also associated with greater instability. It shares the same code CM8. Chaotic altocumulus, which is typically poorly defined with multiple species or transitional forms arranged in several layers, is coded CM9.
Opacity-based varieties: Altocumulus stratiformis has three opacity-based varieties; Translucidus (CM3), perlucidus (CM3 or 7 depending on predominant opacity), and opacus (CM7). Varieties based on opacity are not commonly associated with the species lenticularis, castellanus, or floccus.
Pattern-based varieties: Radiatus (arranged in parallel bands) is sometimes seen with the stratformis and castellanus species. Altocumulus stratiformis radiatus of any opacity is coded CM5 if it is increasing in amount. The duplicatus or undulatus varieties are occasionally seen with the stratiformis and lenticularis species. Altocumulus stratiformis duplicatus is coded CM7 if it is not overridden by another coding of higher importance. Lacunosus is very occasionally associated with altocumulus of the species stratiformis, castellanus, or floccus.
Precipitation-based supplementary feature: Altocumulus often produces virga but usually not precipitation that reaches the ground.
Cloud-based supplementary feature: Mamma caused by localized downdrafts in the cloud layer are occasionally seen with altocumulus. A newly recognized type is the asperitas feature that is characterized by chaotic undulations caused by severe wind shear.
Genitus Mother clouds: Altocumulus stratiformis cumulogenitus or cumulonimbogenitus can form when the middle or upper part of a towering free convective cloud begins to spread horizontally due to a loss of convective lift. It is coded CM6.
Mutatus mother clouds: Altocumulus can form due to the complete transformation of cirrocumulus, altostratus, nimbostratus, or stratocumulus.
Gallery
| Physical sciences | Clouds | Earth science |
166735 | https://en.wikipedia.org/wiki/Altostratus%20cloud | Altostratus cloud | Altostratus is a middle-altitude cloud genus made up of water droplets, ice crystals, or a mixture of the two. Altostratus clouds are formed when large masses of warm, moist air rise, causing water vapor to condense. Altostratus clouds are usually gray or blueish featureless sheets, although some variants have wavy or banded bases. The sun can be seen through thinner altostratus clouds, but thicker layers can be quite opaque.
Altostratus clouds usually predict the arrival of warm fronts. Once altostratus clouds associated with a warm front arrive, continuous rain or snow will usually follow in the next 12 to 24 hours. Although altostratus clouds predict the arrival of warmer, wetter weather, they themselves do not produce significant precipitation. Thunderstorms can be embedded in altostratus clouds, however, bringing showers.
Because altostratus clouds can contain ice crystals, they can produce some optical phenomena like iridescence and coronas.
Description
Altostratus clouds are generally gray or blue-tinged with a largely-uniform blanket-like appearance. They do not have distinct features, and usually do not produce precipitation. The name "altostratus" comes from the conjugation of the Latin words "altum", meaning "high", and "stratus", meaning "flat" or "spread out". Altostratus clouds can produce virga, causing the cloud base to appear hazy. While they do not produce significant precipitation, altostratus clouds can cause light sprinkles or even small rain showers. Consistent rainfall and lowering of the cloud base causes altostratus to become nimbostratus.
Unlike most other types of clouds, altostratus clouds are not subdivided into cloud species due to their largely-featureless appearance. However, they still appear in five varieties: Altostratus duplicatus, opacus, radiatus, translucidus, and undulatus. Altostratus duplicatus is a rare form of altostratus clouds composed of two or more layers of cloud. Translucidus is a translucent form of altostratus clouds, meaning that the sun or moon can be seen through the cloud, whereas the opacus variety is opaque. Radiatus is another rare variety. It has parallel bands of cloud that stretch toward the horizon. The undulatus variety has an wavy appearance—the underside of the cloud appears to rise and fall.
Altostratus and altocumulus clouds, both of which are mid-level clouds, are commonly measured together in cloud cover studies. Together, they cover around 25% of the Earth's surface on average based on CALIPSO satellite data. This constitutes roughly one third of the Earth's total cloud cover. By itself, separated from altocumulus, altostratus covers ~16% of the Earth's surface. Altostratus cloud cover varies seasonally in temperate regions, with significantly less coverage in the summer months as compared to the other seasons. Additionally, altostratus cloud cover varies by latitude, with tropical regions having vastly fewer altostratus clouds when compared to temperate or polar regions. Altostratus and altocumulus cover roughly 22% of the ocean's surface based on surface measurements, with minimal variation based on season.
Altostratus clouds are warmest at the bottom and coldest at the top, with a fairly consistent lapse rate of 5 to 7 °C per kilometer (14 to 20 °F per mile) inside the cloud. The lapse rate is the rate at which the temperature decreases with altitude. Higher lapse rates (i.e. the faster temperature drops with increasing altitude) were associated with colder clouds. The average temperature of altostratus clouds, based on data collected from roughly 45° to 80° latitude, varied from around . Warmer temperatures occurred during summer and colder temperatures during winter.
Inside altostratus clouds, the relative humidity is generally greatest towards the top of the cloud decreasing slowly and roughly linearly towards the bottom. The lowest part of the cloud has the lowest relative humidity. Below the bottom of the cloud, the relative humidity drops rapidly.
Microphysical properties
Altostratus can be composed of water droplets, supercooled water droplets, and ice crystals, but ice crystals make up the vast majority. In some altostratus clouds made of ice crystals, very thin horizontal sheets of water droplets can appear seemingly at random, but they quickly disappear. The sizes of the ice crystals in the cloud tended to increase as altitude decreased. However, close to the bottom of the cloud, the particles decreased in size again. During the sampling of one cloud, the scientists noted a halo while flying near the top of the cloud, which indicated that the ice crystals were hexagonal near the top. However, farther down, the ice crystals became more conglomerated. Mixed-phase (containing both ice and water) altostratus clouds contain a "melt layer", below which the ice crystals tend to melt into water droplets. These water droplets are spheres and thus fall much faster than ice crystals, collecting at the bottom of the cloud.
Formation
Altostratus clouds form when a large mass of warm air rises, causing water vapor in the atmosphere to condense onto nuclei (small dust particles), forming water droplets and ice crystals. These conditions usually happen at the leading edge of a warm front, where cirrostratus clouds thicken and lower until they transition into altostratus clouds. Alternatively, nimbostratus clouds can thin into altostratus. Altostratus can even form from the spreading of the upper anvil cloud or the middle column of a thunderstorm.
Altostratus clouds are mid-level clouds that form from above sea level in polar regions. In temperate regions, the ceiling increases drastically, allowing altostratus clouds to form between . In tropical regions, altostratus can reach even higher, forming from . They can range from in thickness and can cover hundreds of kilometers of the Earth's surface.
Use in forecasting
Altostratus clouds tend to form ahead of warm fronts or occluded fronts and herald their arrival. These warm fronts bring warmer air into the region. Occluded fronts form when a faster-moving cold front catches up to a warm front, and the temperature after the frontal system passes may rise or fall. As the frontal system approaches, cirrostratus clouds will thicken into altostratus clouds, which then gradually thicken further into nimbostratus clouds. If the frontal system is occluded, cumulonimbus clouds may also be present. Once the altostratus clouds have arrived, rain or snow will usually follow in the next 12 to 24 hours.
Instability in the atmosphere can embed thunderstorms in an altostratus cloud, although altostratus clouds themselves do not produce storms.
Effects on climate
Globally, clouds reflect around 50 watts per square meter of short-wave solar radiation back into space, cooling the Earth by around , an effect largely caused by stratocumulus clouds. However, at the same time, they reflect around 30 watts per square meter of long-wave (infrared) black body radiation emitted by the Earth back to Earth's surface, heating the Earth by around —a process called the greenhouse effect. Cirrus and altostratus clouds are the top two sources of this heating effect. This combination of heating and cooling sums out to a net loss of 20 watts per square meter globally, cooling the Earth by roughly .
Altostratus clouds are the only cloud genus besides cirrus clouds to exhibit a net global heating effect on Earth and its atmosphere; however, cirrus have a heating effect that is four times as potent as altostratus (2 watts per square meter versus only 0.5 watts per square meter).
Optical phenomena
Altostratus clouds can produce bright halos when viewed from the air, but not when viewed from the ground. Halos can take the appearance of rings, arcs, or spots of white or multicolored light and are formed by the reflection and refraction of sunlight or moonlight shining through ice crystals in the cloud. Light diffraction through altostratus clouds can also produce coronas, which are small, concentric pastel-colored rings of light around the sun or moon. They can also be iridescent, with often-parallel bands of bright color projected on a cloud. Unlike the halos, the coronas and iridescence can be seen from Earth's surface.
Relation to other clouds
Altostratus and altocumulus clouds are the two genera of mid-level clouds that usually form between . These are given the prefix "alto-". These clouds are formed from ice crystals, supercooled water droplets, or liquid water droplets.
Above the mid-level clouds are three different genera of high-level clouds, cirrus, cirrocumulus, and cirrostratus, all of which are given the prefix "cirro-". High-level clouds usually form above . Cirrocumulus and cirrostratus are sometimes informally referred to as cirriform clouds because of their frequent association with cirrus.
Below the mid-level clouds are the low-level clouds, which usually form below and do not have a prefix. The two genera that are strictly low-level are stratus, and stratocumulus. These clouds are composed of water droplets, except during winter when they are formed of supercooled water droplets or ice crystals if the temperature at cloud level is below freezing. Three additional genera usually form in the low altitude range, but may be based at higher levels under conditions of very low humidity. They are the genera cumulus, and cumulonimbus, and nimbostratus. These are sometimes classified separately as clouds of vertical development, especially when their tops are high enough to be composed of supercooled water droplets or ice crystals.
Cirrostratus
Cirrostratus clouds can appear as a smooth veil in the sky or as a striated sheet. They are sometimes similar to altostratus and are distinguishable from the latter because the sun or moon is always clearly visible through transparent cirrostratus, in contrast to altostratus which tends to be opaque or translucent. Cirrostratus come in two species, fibratus and nebulosus. The ice crystals in these clouds vary depending upon the height in the cloud. Towards the bottom, at temperatures of around , the crystals tend to be long, solid, hexagonal columns. Towards the top of the cloud, at temperatures of around , the predominant crystal types are thick, hexagonal plates and short, solid, hexagonal columns. These clouds commonly produce halos, and sometimes the halo is the only indication that such clouds are present. They are formed by warm, moist air being lifted slowly to a very high altitude. When a warm front approaches, cirrostratus clouds become thicker and descend forming altostratus clouds, and rain usually begins 12 to 24 hours later.
Altocumulus
Altocumulus clouds are small patches or heaps of white or light gray cloud. Like altostratus, altocumulus are composed of a mixture of water droplets, supercooled water droplets, and ice crystals. Although altocumulus clouds are mid-level clouds that form at roughly the same altitude as altostratus clouds, their formation methods are completely different. Altocumulus forms from convective (rising) processes, whereas altostratus is usually formed by descending and thickening cirrostratus.
Stratus
Stratus are low-level clouds that are usually visually similar to altostratus. Stratus comes in two species: nebulosus, a largely-featureless flat gray cloud sheet, and fractus, shattered fragments of cloud often called "scud". Opaque varieties of altostratus and stratus nebulosus clouds can be virtually indistinguishable from each other to the naked eye, to the point that the World Meteorological Organization suggests that one of the few ways to distinguish between these clouds is to check what types of clouds came before them. Altostratus clouds, because they tend to form from warm fronts, are usually preceded by high-level cirriform clouds. Stratus clouds tend to form by cooling air masses, often at night, and thus are not usually preceded by other types of clouds.
Nimbostratus
Nimbostratus are low-level (sometimes classified as vertical) rain-bearing stratus clouds. Unlike the sprinkles or light drizzles that altostratus or stratus can produce, nimbostratus produces heavy, continuous rain or snow. These clouds are thick and dark enough to entirely blot out the sun. Nimbostratus has no species or varieties. Like altostratus, nimbostratus clouds can be made of ice crystals, supercooled water droplets, or water droplets.
| Physical sciences | Clouds | Earth science |
166736 | https://en.wikipedia.org/wiki/Stratocumulus%20cloud | Stratocumulus cloud | }}
A stratocumulus cloud, occasionally called a cumulostratus, belongs to a genus-type of clouds characterized by large dark, rounded masses, usually in groups, lines, or waves, the individual elements being larger than those in altocumulus, and the whole being at a lower height, usually below . Weak convective currents create shallow cloud layers (see also: sea of clouds) because of drier, stable air above preventing continued vertical development. Historically, in English, this type of cloud has been referred to as a twain cloud for being a combination of two types of clouds.
Description
Stratocumulus clouds are rounded clumps or patches of white to dark gray clouds that normally form in groups. The individual cloud elements, which cover more than 5 degrees of arc each, can connect with each other and are sometimes arranged in a regular pattern.
Occurrence
Vast areas of subtropical and polar oceans are covered with massive sheets of stratocumulus. These may organize into distinctive patterns which are currently under active study. In subtropics, they cover the edges of the horse latitude climatological highs, and reduce the amount of solar energy absorbed in the ocean. When these drift over land the summer heat or winter cold is reduced. 'Dull weather' is a common expression incorporated with overcast stratocumulus days, which usually occur either in a warm sector between a warm and cold front in a depression, or in an area of high pressure, in the latter case, sometimes persisting over a specific area for several days. If the air over land is moist and hot enough, stratocumulus may develop to various cumulus clouds, or, more commonly, the sheet of stratocumulus may become thick enough to produce some light rain. On drier areas they quickly dissipate over land, resembling cumulus humilis. This often occurs in late morning in areas under anticyclonic weather, the stratocumulus breaking up under the sun's heat and often reforming again by evening as the heat of the sun decreases again.
Precipitation
Most often, stratocumulus produce no precipitation, and when they do, it is generally only light rain or snow. However, these clouds are often seen at either the front or tail end of worse weather, so they may indicate storms to come, in the form of thunderheads or gusty winds. They are also often seen underneath the cirrostratus and altostratus sheets that often precede a warm front, as these higher clouds decrease the sun's heat and therefore convection, causing any cumulus clouds to spread out into stratocumulus clouds.
Comparison with altocumulus
Stratocumulus clouds are similar in appearance to altocumulus and can be mistaken for such. A simple test to distinguish these is to compare the size of individual masses or rolls: when pointing one's hand in the direction of the cloud, if the cloud is about the size of the thumb, it is altocumulus; if it is the size of one's fist, it is stratocumulus. This often does not apply when stratocumulus is of a broken, fractus form, when it may appear as small as altocumulus.
Optical effects
Stratocumulus clouds are the main type of cloud that can produce crepuscular rays. Thin stratocumulus clouds are also often the cause of corona effects around the Moon at night. All stratocumulus subtypes are coded CL5 except when formed from free convective mother clouds (CL4) or when formed separately from co-existing (CL8).
Formation
Stratocumulus clouds usually form from the rising and breakup of a stratus cloud. They can also form from altostratus and nimbostratus clouds, either as evaporating precipitation condenses into a cloud or as the nimbostratus cloud itself thins and breaks up. If a cumulus cloud becomes flattened (for example, by wind shear or temperature inversion), it too can become a stratocumulus cloud.
Species
Stratocumulus Stratiformis are extensive flat but slightly lumpy sheets that show only minimal convective activity.
Stratocumulus Lenticularis are separate flat elongated seed-shaped clouds. They are typical for polar countries or warmer climate during winter seasons. They also can be formed by winds passing hills or mountains, such as Foehn winds, and in this case they can be very regularly shaped.
Stratocumulus Castellanus have stronger convective activity due to the presence of increasingly unstable air. They are distinct from other stratocumulus by puffy tower-like formations atop the cloud layer. They look like cumulus congestus, but can be easily confused: "towers" of cumulus congestus grow above separate clouds, whereas in the case of stratocumulus castellanus, there is always a more or less defined layer of clouds. Stratocumulus castellanus may develop into cumulus congestus (and even further into cumulonimbus) under auspicious conditions. Any showers from stratocumulus castellanus are not usually as heavy as those from cumulus congestus.
Opacity-based varieties
Stratocumulus Opacus is a dark layer of clouds covering entire sky without any break. However, the cloud sheet is not completely uniform, so that separate cloud bases still can be seen. This is the main precipitating type, however any rain is usually light. If the cloud layer becomes grayer to the point when individual clouds cannot be distinguished, stratocumulus turn into stratus clouds.
Stratocumulus Perlucidus is a layer of stratocumulus clouds with small spaces, appearing in irregular pattern, through which clear sky or higher clouds can be seen.
Stratocumulus Translucidus consist of separate groups of stratocumulus clouds, with a clear sky (or higher clouds) visible between them. No precipitation in most cases.
Pattern-based varieties
Stratocumulus Undulatus clouds appear as nearly parallel waves, rolls or separate elongated clouds, without significant vertical development.
Stratocumulus Radiatus clouds appear as the same as stratocumulus undulatus, but stratocumulus undulatus move perpendicular to the wind shear, while stratocumulus radiatus move parallel to the wind shear.
Stratocumulus Duplicatus clouds appear as stratocumulus clouds with two or more layers or sheets. Stratocumulus duplicatus is common on species lenticularis or lenticular cloud.
Stratocumulus Lacunosus clouds are very uncommon. They only occur when there are localized downdrafts striking through the stratocumuliform cloud.
Supplementary feature
Stratocumulus Mamma is a type of mammatus cloud.
Stratocumulus Asperitas is a rare, newly recognized supplementary feature that presents itself as chaotic, wavy undulations appearing in the base of a stratocumulus cloud cover. It is thought these clouds are formed by severe wind shear.
Stratocumulus Fluctus is also a rare, newly recognized supplementary feature in which short-lived "sea waves" form on top of a stratocumulus cloud, they are caused by wind speed and direction differences directly under and over the cloud.
Precipitation-based supplementary features
Stratocumulus Virga is a form of precipitation that evaporates in mid-air and doesn't reach the ground.
Stratocumulus Praecipitatio is a form of precipitation that reaches the ground as light rain or snow.
Mother clouds
Stratocumulus Cumulomutatus the specific type of stratocumulus clouds, are flat and elongated. They form in the evening, when updrafts caused by convection decrease making cumulus clouds lose vertical development and spread horizontally. They also can occur under altostratus cloud preceding a warm or occluded front, when cumulus usually lose vertical development as the sun's heat decreases. Like all other forms of stratocumulus apart from castellanus, they are also often found in anticyclones.
Stratocumulus Cumulogenitus out of cumulus or cumulonimbus clouds, disrupted by decreasing convection. During formation period, puffy tops of cumulus clouds can protrude from stratocumulus cumulogenitus for a relatively long time until they completely spread in horizontal direction. Stratocumulus cumulogenitus appear as lengthy sheet or as group of separate elongated cloud rolls or waves.
Hypothetical break-up
| Physical sciences | Clouds | Earth science |
166737 | https://en.wikipedia.org/wiki/Nimbostratus%20cloud | Nimbostratus cloud | A nimbostratus cloud is a multilevel, amorphous, nearly uniform, and often dark-grey cloud that usually produces continuous rain, snow, or sleet, but no lightning or thunder.
Although it is usually a low-based stratiform cloud, it actually forms most commonly in the middle level of the troposphere and then spreads vertically into the low and high levels. Nimbostratus usually produces precipitation over a wide area.
The prefix nimbo- comes from the Latin word , which denotes "dark cloud".
Downward-growing nimbostratus can have the same vertical extent as most large upward-growing cumulus, but its horizontal expanse tends to be even greater.
Appearance
Nimbostratus has a diffuse cloud base generally found anywhere from near surface in the low levels to about .in the middle level of the troposphere. Although usually dark at its base, it often appears illuminated from within to a surface observer. Though found worldwide, nimbostratus occurs more commonly in the middle latitudes. It is coded CM2 on the SYNOP report.
Formation
Nimbostratus occurs along a warm front or occluded front where the slowly rising warm air mass creates nimbostratus along with shallower stratus clouds producing less rain, these clouds being preceded by higher-level clouds such as cirrostratus and altostratus. Often, when an altostratus cloud thickens and descends into lower altitudes, it will become nimbostratus.
Nimbostratus, unlike cumulonimbus, is not associated with thunderstorms, however at an unusually unstable warm front caused as a result of the advancing warm air being hot, humid and unstable, cumulonimbus clouds may be embedded within the usual nimbostratus. Lightning from an embedded cumulonimbus cloud may interact with the nimbostratus but only in the immediate area around it. In this situation with lightning and rain occurring it would be hard to tell which type of cloud was producing the rain from the ground, however cumulonimbus tend to produce larger droplets and more intense downpours. The occurrence of cumulonimbus and nimbostratus together is uncommon, and usually only nimbostratus is found at a warm front and sometimes in cold front.
Forecast
Nimbostratus is generally a sign of an approaching warm or occluded front producing steady moderate precipitation, as opposed to the shorter period of typically heavier precipitation released by a cold-frontal cumulonimbus cloud. Precipitation may last for several days, depending on the speed of the frontal system. A nimbostratus virga cloud is the same as normal nimbostratus, but the rain or snow falls as virga which doesn't reach the ground. Stratus or stratocumulus usually replace the nimbostratus after the passage of the warm or occluded front.
Origin of name
Under Luke Howard's first systematized study of clouds, carried out in France in 1802, three general cloud forms were established based on appearance and characteristics of formation: cirriform, cumuliform and stratiform. These were further divided into upper and lower types depending on altitude. In addition to these three main types, Howard added two names to designate multiple cloud types joined together: cumulostratus, a blending of cumulus clouds and stratus layers, and nimbus, a complex blending of cirriform, cumuliform, and stratiform clouds with sufficient vertical development to produce significant precipitation.
Later, in the 20th century, an IMC commission for the study of clouds put forward a refined and more restricted definition of the genus nimbus, effectively reclassifying it as a stratiform cloud type. It was then renamed nimbostratus, and published with the new name in the 1932 edition of the International Atlas of Clouds and of States of the Sky. This left cumulonimbus as the only nimboform type as indicated by its root name.
Subtypes and derivative types
Species and varieties: Nimbostratus is very thick, opaque, and featureless, so this genus type is not subdivided into species or varieties.
Precipitation-based supplementary features: Nimbostratus is a major precipitation cloud and produces the virga or precipitation features. The latter can achieve heavy intensity due to the cloud's vertical depth.
Accessory cloud: Nimbostratus pannus is an accessory cloud of nimbostratus that forms as a ragged layer in precipitation below the main cloud deck. Pannus is coded CL7.
Genitus mother clouds: This genus type can form from cumulus and cumulonimbus.
Mutatus mother clouds: Nimbostratus can form due to the complete transformation of altocumulus, altostratus and stratocumulus.
Relation to other clouds
Multi-level nimbostratus is physically related to other stratiform genus-types by way of being non-convective in nature. However, the other sheet-like clouds usually each occupy only one or two levels at the same time. Stratus clouds are low-level and form from near ground level to at all latitudes. In the middle level are the altostratus clouds that form from to in polar areas, in temperate areas, and in tropical areas. Although altostratus forms mostly in the middle level of the troposphere, strong frontal lift can push it into the lower part of the high-level. The main high-level stratiform cloud is cirrostratus which is composed of ice crystals that often produce halo effects around the sun. Cirrostratus forms at altitudes of in high latitudes, in temperate latitudes, and in low, tropical latitudes. Of the non-stratiform clouds, cumulonimbus and cumulus congestus are the most closely related to nimbostratus because of their vertical extent and ability to produce moderate to heavy precipitation. The remaining cumuliform (cumulus) and stratocumuliform (stratocumulus, altocumulus, and cirrocumulus) clouds have the least in common with nimbostratus.
| Physical sciences | Clouds | Earth science |
166740 | https://en.wikipedia.org/wiki/Cirrocumulus%20cloud | Cirrocumulus cloud | Cirrocumulus is one of the three main genus types of high-altitude tropospheric clouds, the other two being cirrus and cirrostratus. They usually occur at an altitude of , however they can occur as low as in the arctic and weather reporting standards such as the Canadian MANOBS suggests heights of in summer and in winter. Like lower-altitude cumuliform and stratocumuliform clouds, cirrocumulus signifies convection. Unlike other high-altitude tropospheric clouds like cirrus and cirrostratus, cirrocumulus includes a small amount of liquid water droplets, although these are in a supercooled state. Ice crystals are the predominant component, and typically, the ice crystals cause the supercooled water drops in the cloud to rapidly freeze, transforming the cirrocumulus into cirrostratus. This process can also produce precipitation in the form of a virga consisting of ice or snow. Thus, cirrocumulus clouds are usually short-lived. They usually only form as part of a short-lived transitional phase within an area of cirrus clouds and can also form briefly as a result of the breaking up of part of a cumulonimbus anvil.
Properly, the term cirrocumulus refers to each cloud, but is typically also used to refer to an entire patch of cirrocumulus. When used in this way, each cirrocumulus element is referred to as a "cloudlet".
Appearance
Cirrocumulus is a cloud of the stratocumuliform physical category that shows both stratiform and cumuliform characteristics and typically appears as white, patchy sheets with ripples or tufts without gray shading. Each cloudlet appears no larger than a finger held at arm's length. These often are organized in rows like other cumuliform and stratocumuliform clouds, but since they are so small, cirrocumulus patches take on a finer appearance, sometimes also referred to colloquially as "herringbone" or as a "mackerel sky". Cirrocumulus is coded CH9 for the main genus-type and all subforms.
Cirrocumulus is distinguished from altocumulus in several ways, although the two stratocumuliform genus types can occasionally occur together with no clear demarcation between them. Cirrocumulus generally occur at higher altitudes than altocumulus, thus the "cloudlets" appear smaller, as they are more distant from observation at ground level. They are also colder. Cirrocumulus clouds never cast self-shadows and are translucent to a certain degree. They are also typically found amongst other cirrus clouds in the sky and are usually themselves seen to be transforming into these other types of cirrus. This often occurs at the leading edge of a warm front, where many types of cirriform clouds can be present.
Cirrocumulus clouds tend to reflect the red and yellow colours during a sunset and sunrise, so they have been referred to as "one of the most beautiful clouds". This occurs because they reflect the unscattered rays of light from the early morning or evening sun, and those rays are yellow, orange, red, and sometimes purple.
Forecasting
Cirrocumulus usually only forms in patches. If it forms in patches with cirrus or cirrostratus and the clouds spread across the sky, it usually means rain in 8–10 hours (can be more if the front is slow-moving). Only small patches of cirrocumulus and perhaps some wisps of cirrus usually mean a continuation of good weather (although this may also be seen in conjunction with showers and thunderstorms). If it is seen after rain, it usually means improving the weather.
Subtypes
Species: Cirrocumulus stratiformis (Cc str) is one of four species and appears in the form of relatively flat stratocumuliform sheets or patches. The species Cirrocumulus lenticularis (Cc len) takes its name from the lens-shaped structure of this cloud which is tapered at each end. Cirrocumulus castellanus (Cc cas) has cumuliform buildups that give the cloud a partly or mainly turreted appearance. When the cumuliform parts have more of a tufted appearance, it is given the species name Cirrocumulus floccus (Cc flo). Cirrocumulus with castellanus buildups can show some vertical extent, but are not usually classified as vertical or multi-étage clouds.
Varieties: This genus type is always translucent and so has no opacity-based varieties. However, like cirrus, certain cirrocumulus species can sometimes be divided into pattern-based varieties. The undulatus variety has a wavy undulating base and is seen mostly with the stratiformis and lenticularis species types. The lacunosus variety contains circular holes caused by downdrafts in the cloud and is associated mainly with the species stratiformis, castellanus, and floccus.
Precipitation-based supplementary feature: Cirrocumulus occasionally produces virga, precipitation that evaporates before reaching the ground.
Cloud-based supplementary feature: Mamma in the form of downward forming bubbles is infrequently seen as a cloud-based supplementary feature.
Mother clouds: This genus type has no recognized genitus mother clouds. However cirrocumulus stratiformis cirromutatus or cirrostratomutatus can result from sheets or filaments of high cloud taking on a stratocumuli form structure as a result of high altitude convection. A high layer of white or light grey altocumulus of a particular species can thin out into pure white cirrocumulus altocumulomutatus of the same species.
| Physical sciences | Clouds | Earth science |
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