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55951 | https://en.wikipedia.org/wiki/Instant%20messaging | Instant messaging | Instant messaging (IM) technology is a type of synchronous computer-mediated communication involving the immediate (real-time) transmission of messages between two or more parties over the Internet or another computer network. Originally involving simple text message exchanges, modern IM applications and services (also called "social messengers", "messaging apps", "chat apps" or "chat clients") tend to also feature the exchange of multimedia, emojis, file transfer, VoIP (voice calling), and video chat capabilities.
Instant messaging systems facilitate connections between specified known users (often using a contact list also known as a "buddy list" or "friend list") or in chat rooms, and can be standalone apps or integrated into a wider social media platform, or in a website where it can, for instance, be used for conversational commerce. Originally the term "instant messaging" was distinguished from "text messaging" by being run on a computer network instead of a cellular/mobile network, being able to write longer messages, real-time communication, presence ("status"), and being free (only cost of access instead of per SMS message sent).
Instant messaging was pioneered in the early Internet era; the IRC protocol was the earliest to achieve wide adoption. Later in the 1990s, ICQ was among the first closed and commercialized instant messengers, and several rival services appeared afterwards as it became a popular use of the Internet. Beginning with its first introduction in 2005, BlackBerry Messenger became the first popular example of mobile-based IM, combining features of traditional IM and mobile SMS. Instant messaging remains very popular today; IM apps are the most widely used smartphone apps: in 2018 for instance there were 980 million monthly active users of WeChat and 1.3 billion monthly users of WhatsApp, the largest IM network.
Overview
Instant messaging (IM), sometimes also called "messaging" or "texting", consists of computer-based human communication between two users (private messaging) or more (chat room or "group") in real-time, allowing immediate receipt of acknowledgment or reply. This is in direct contrast to email, where conversations are not in real-time, and the perceived quasi-synchrony of the communications by the users (although many systems allow users to send offline messages that the other user receives when logging in).
Earlier IM networks were limited to text-based communication, not dissimilar to mobile text messaging. As technology has moved forward, IM has expanded to include voice calling using a microphone, videotelephony using webcams, file transfer, location sharing, image and video transfer, voice notes, and other features.
IM is conducted over the Internet or other types of networks (see also LAN messenger). Depending on the IM protocol, the technical architecture can be peer-to-peer (direct point-to-point transmission) or client–server (when all clients have to first connect to the central server). Primary IM services are controlled by their corresponding companies and usually follow the client-server model.
The term "Instant Messenger" is a service mark of Time Warner and may not be used in software not affiliated with AOL in the United States. For this reason, in April 2007, the instant messaging client formerly named Gaim (or gaim) announced that they would be renamed "Pidgin".
Clients
Modern IM services generally provide their own client, either a separately installed software or a browser-based client. They are normally centralised networks run by the servers of the platform's operators, unlike peer-to-peer protocols like XMPP. These usually only work within the same IM network, although some allow limited function with other services (see #Interoperability). Third-party client software applications exist that will connect with most of the major IM services. There is the class of instant messengers that uses the serverless model, which doesn't require servers, and the IM network consists only of clients. There are several serverless messengers: RetroShare, Tox, Bitmessage, Ricochet, Ring. | Technology | Internet | null |
55955 | https://en.wikipedia.org/wiki/Version%20control | Version control | Version control (also known as revision control, source control, and source code management) is the software engineering practice of controlling, organizing, and tracking different versions in history of computer files; primarily source code text files, but generally any type of file.
Version control is a component of software configuration management.
A version control system is a software tool that automates version control. Alternatively, version control is embedded as a feature of some systems such as word processors, spreadsheets, collaborative web docs, and content management systems, e.g., Wikipedia's page history.
Version control includes viewing old versions and enables reverting a file to a previous version.
Overview
As teams develop software, it is common for multiple versions of the same software to be deployed in different sites and for the developers to work simultaneously on updates. Bugs or features of the software are often only present in certain versions (because of the fixing of some problems and the introduction of others as the program develops). Therefore, for the purposes of locating and fixing bugs, it is vitally important to be able to retrieve and run different versions of the software to determine in which version(s) the problem occurs. It may also be necessary to develop two versions of the software concurrently: for instance, where one version has bugs fixed, but no new features (branch), while the other version is where new features are worked on (trunk).
At the simplest level, developers could simply retain multiple copies of the different versions of the program, and label them appropriately. This simple approach has been used in many large software projects. While this method can work, it is inefficient as many near-identical copies of the program have to be maintained. This requires a lot of self-discipline on the part of developers and often leads to mistakes. Since the code base is the same, it also requires granting read-write-execute permission to a set of developers, and this adds the pressure of someone managing permissions so that the code base is not compromised, which adds more complexity. Consequently, systems to automate some or all of the revision control process have been developed. This ensures that the majority of management of version control steps is hidden behind the scenes.
Moreover, in software development, legal and business practice, and other environments, it has become increasingly common for a single document or snippet of code to be edited by a team, the members of which may be geographically dispersed and may pursue different and even contrary interests. Sophisticated revision control that tracks and accounts for ownership of changes to documents and code may be extremely helpful or even indispensable in such situations.
Revision control may also track changes to configuration files, such as those typically stored in /etc or /usr/local/etc on Unix systems. This gives system administrators another way to easily track changes made and a way to roll back to earlier versions should the need arise.
Many version control systems identify the version of a file as a number or letter, called the version number, version, revision number, revision, or revision level. For example, the first version of a file might be version 1. When the file is changed the next version is 2. Each version is associated with a timestamp and the person making the change. Revisions can be compared, restored, and, with some types of files, merged.
History
IBM's OS/360 IEBUPDTE software update tool dates back to 1962, arguably a precursor to version control system tools. Two source management and version control packages that were heavily used by IBM 360/370 installations were The Librarian and Panvalet.
A full system designed for source code control was started in 1972, Source Code Control System for the same system (OS/360). Source Code Control System's introduction, having been published on December 4, 1975, historically implied it was the first deliberate revision control system. RCS followed just after, with its networked version Concurrent Versions System. The next generation after Concurrent Versions System was dominated by Subversion, followed by the rise of distributed revision control tools such as Git.
Structure
Revision control manages changes to a set of data over time. These changes can be structured in various ways.
Often the data is thought of as a collection of many individual items, such as files or documents, and changes to individual files are tracked. This accords with intuitions about separate files but causes problems when identity changes, such as during renaming, splitting or merging of files. Accordingly, some systems such as Git, instead consider changes to the data as a whole, which is less intuitive for simple changes but simplifies more complex changes.
When data that is under revision control is modified, after being retrieved by checking out, this is not in general immediately reflected in the revision control system (in the repository), but must instead be checked in or committed. A copy outside revision control is known as a "working copy". As a simple example, when editing a computer file, the data stored in memory by the editing program is the working copy, which is committed by saving. Concretely, one may print out a document, edit it by hand, and only later manually input the changes into a computer and save it. For source code control, the working copy is instead a copy of all files in a particular revision, generally stored locally on the developer's computer; in this case saving the file only changes the working copy, and checking into the repository is a separate step.
If multiple people are working on a single data set or document, they are implicitly creating branches of the data (in their working copies), and thus issues of merging arise, as discussed below. For simple collaborative document editing, this can be prevented by using file locking or simply avoiding working on the same document that someone else is working on.
Revision control systems are often centralized, with a single authoritative data store, the repository, and check-outs and check-ins done with reference to this central repository. Alternatively, in distributed revision control, no single repository is authoritative, and data can be checked out and checked into any repository. When checking into a different repository, this is interpreted as a merge or patch.
Graph structure
In terms of graph theory, revisions are generally thought of as a line of development (the trunk) with branches off of this, forming a directed tree, visualized as one or more parallel lines of development (the "mainlines" of the branches) branching off a trunk. In reality the structure is more complicated, forming a directed acyclic graph, but for many purposes "tree with merges" is an adequate approximation.
Revisions occur in sequence over time, and thus can be arranged in order, either by revision number or timestamp. Revisions are based on past revisions, though it is possible to largely or completely replace an earlier revision, such as "delete all existing text, insert new text". In the simplest case, with no branching or undoing, each revision is based on its immediate predecessor alone, and they form a simple line, with a single latest version, the "HEAD" revision or tip. In graph theory terms, drawing each revision as a point and each "derived revision" relationship as an arrow (conventionally pointing from older to newer, in the same direction as time), this is a linear graph. If there is branching, so multiple future revisions are based on a past revision, or undoing, so a revision can depend on a revision older than its immediate predecessor, then the resulting graph is instead a directed tree (each node can have more than one child), and has multiple tips, corresponding to the revisions without children ("latest revision on each branch"). In principle the resulting tree need not have a preferred tip ("main" latest revision) – just various different revisions – but in practice one tip is generally identified as HEAD. When a new revision is based on HEAD, it is either identified as the new HEAD, or considered a new branch. The list of revisions from the start to HEAD (in graph theory terms, the unique path in the tree, which forms a linear graph as before) is the trunk or mainline. Conversely, when a revision can be based on more than one previous revision (when a node can have more than one parent), the resulting process is called a merge, and is one of the most complex aspects of revision control. This most often occurs when changes occur in multiple branches (most often two, but more are possible), which are then merged into a single branch incorporating both changes. If these changes overlap, it may be difficult or impossible to merge, and require manual intervention or rewriting.
In the presence of merges, the resulting graph is no longer a tree, as nodes can have multiple parents, but is instead a rooted directed acyclic graph (DAG). The graph is acyclic since parents are always backwards in time, and rooted because there is an oldest version. Assuming there is a trunk, merges from branches can be considered as "external" to the tree – the changes in the branch are packaged up as a patch, which is applied to HEAD (of the trunk), creating a new revision without any explicit reference to the branch, and preserving the tree structure. Thus, while the actual relations between versions form a DAG, this can be considered a tree plus merges, and the trunk itself is a line.
In distributed revision control, in the presence of multiple repositories these may be based on a single original version (a root of the tree), but there need not be an original root - instead there can be a separate root (oldest revision) for each repository. This can happen, for example, if two people start working on a project separately. Similarly, in the presence of multiple data sets (multiple projects) that exchange data or merge, there is no single root, though for simplicity one may think of one project as primary and the other as secondary, merged into the first with or without its own revision history.
Specialized strategies
Engineering revision control developed from formalized processes based on tracking revisions of early blueprints or bluelines. This system of control implicitly allowed returning to an earlier state of the design, for cases in which an engineering dead-end was reached in the development of the design. A revision table was used to keep track of the changes made. Additionally, the modified areas of the drawing were highlighted using revision clouds.
In Business and Law
Version control is widespread in business and law. Indeed, "contract redline" and "legal blackline" are some of the earliest forms of revision control, and are still employed in business and law with varying degrees of sophistication. The most sophisticated techniques are beginning to be used for the electronic tracking of changes to CAD files (see product data management), supplanting the "manual" electronic implementation of traditional revision control.
Source-management models
Traditional revision control systems use a centralized model where all the revision control functions take place on a shared server. If two developers try to change the same file at the same time, without some method of managing access the developers may end up overwriting each other's work. Centralized revision control systems solve this problem in one of two different "source management models": file locking and version merging.
Atomic operations
An operation is atomic if the system is left in a consistent state even if the operation is interrupted. The commit operation is usually the most critical in this sense. Commits tell the revision control system to make a group of changes final, and available to all users. Not all revision control systems have atomic commits; Concurrent Versions System lacks this feature.
File locking
The simplest method of preventing "concurrent access" problems involves locking files so that only one developer at a time has write access to the central "repository" copies of those files. Once one developer "checks out" a file, others can read that file, but no one else may change that file until that developer "checks in" the updated version (or cancels the checkout).
File locking has both merits and drawbacks. It can provide some protection against difficult merge conflicts when a user is making radical changes to many sections of a large file (or group of files). If the files are left exclusively locked for too long, other developers may be tempted to bypass the revision control software and change the files locally, forcing a difficult manual merge when the other changes are finally checked in. In a large organization, files can be left "checked out" and locked and forgotten about as developers move between projects - these tools may or may not make it easy to see who has a file checked out.
Version merging
Most version control systems allow multiple developers to edit the same file at the same time. The first developer to "check in" changes to the central repository always succeeds. The system may provide facilities to merge further changes into the central repository, and preserve the changes from the first developer when other developers check in.
Merging two files can be a very delicate operation, and usually possible only if the data structure is simple, as in text files. The result of a merge of two image files might not result in an image file at all. The second developer checking in the code will need to take care with the merge, to make sure that the changes are compatible and that the merge operation does not introduce its own logic errors within the files. These problems limit the availability of automatic or semi-automatic merge operations mainly to simple text-based documents, unless a specific merge plugin is available for the file types.
The concept of a reserved edit can provide an optional means to explicitly lock a file for exclusive write access, even when a merging capability exists.
Baselines, labels and tags
Most revision control tools will use only one of these similar terms (baseline, label, tag) to refer to the action of identifying a snapshot ("label the project") or the record of the snapshot ("try it with baseline X"). Typically only one of the terms baseline, label, or tag is used in documentation or discussion; they can be considered synonyms.
In most projects, some snapshots are more significant than others, such as those used to indicate published releases, branches, or milestones.
When both the term baseline and either of label or tag are used together in the same context, label and tag usually refer to the mechanism within the tool of identifying or making the record of the snapshot, and baseline indicates the increased significance of any given label or tag.
Most formal discussion of configuration management uses the term baseline.
Distributed revision control
Distributed revision control systems (DRCS) take a peer-to-peer approach, as opposed to the client–server approach of centralized systems. Rather than a single, central repository on which clients synchronize, each peer's working copy of the codebase is a bona-fide repository.
Distributed revision control conducts synchronization by exchanging patches (change-sets) from peer to peer. This results in some important differences from a centralized system:
No canonical, reference copy of the codebase exists by default; only working copies.
Common operations (such as commits, viewing history, and reverting changes) are fast, because there is no need to communicate with a central server.
Rather, communication is only necessary when pushing or pulling changes to or from other peers.
Each working copy effectively functions as a remote backup of the codebase and of its change-history, providing inherent protection against data loss.
Best practices
Following best practices is necessary to obtain the full benefits of version control. Best practice may vary by version control tool and the field to which version control is applied. The generally accepted best practices in software development include: making incremental, small, changes; making commits which involve only one task or fix -- a corollary to this is to commit only code which works and does not knowingly break existing functionality; utilizing branching to complete functionality before release; writing clear and descriptive commit messages, make what why and how clear in either the commit description or the code; and using a consistent branching strategy. Other best software development practices such as code review and automated regression testing may assist in the following of version control best practices.
Costs and benefits
Costs and benefits will vary dependent upon the version control tool chosen and the field in which it is applied. This section speaks to the field of software development, where version control is widely applied.
Costs
In addition to the costs of licensing the version control software, using version control requires time and effort. The concepts underlying version control must be understood and the technical particulars required to operate the version control software chosen must be learned. Version control best practices must be learned and integrated into the organization's existing software development practices. Management effort may be required to maintain the discipline needed to follow best practices in order to obtain useful benefit.
Benefits
Allows for reverting changes
A core benefit is the ability to keep history and revert changes, allowing the developer to easily undo changes. This gives the developer more opportunity to experiment, eliminating the fear of breaking existing code.
Branching simplifies deployment, maintenance and development
Branching assists with deployment. Branching and merging, the production, packaging, and labeling of source code patches and the easy application of patches to code bases, simplifies the maintenance and concurrent development of the multiple code bases associated with the various stages of the deployment process; development, testing, staging, production, etc.
Damage mitigation, accountability and process and design improvement
There can be damage mitigation, accountability, process and design improvement, and other benefits associated with the record keeping provided by version control, the tracking of who did what, when, why, and how.
When bugs arise, knowing what was done when helps with damage mitigation and recovery by assisting in the identification of what problems exist, how long they have existed, and determining problem scope and solutions. Previous versions can be installed and tested to verify conclusions reached by examination of code and commit messages.
Simplifies debugging
Version control can greatly simplify debugging. The application of a test case to multiple versions can quickly identify the change which introduced a bug. The developer need not be familiar with the entire code base and can focus instead on the code that introduced the problem.
Improves collaboration and communication
Version control enhances collaboration in multiple ways. Since version control can identify conflicting changes, i.e. incompatible changes made to the same lines of code, there is less need for coordination among developers.
The packaging of commits, branches, and all the associated commit messages and version labels, improves communication between developers, both in the moment and over time. Better communication, whether instant or deferred, can improve the code review process, the testing process, and other critical aspects of the software development process.
Integration
Some of the more advanced revision-control tools offer many other facilities, allowing deeper integration with other tools and software-engineering processes.
Integrated development environment
Plugins are often available for IDEs such as Oracle JDeveloper, IntelliJ IDEA, Eclipse, Visual Studio, Delphi, NetBeans IDE, Xcode, and GNU Emacs (via vc.el). Advanced research prototypes generate appropriate commit messages.
Common terminology
Terminology can vary from system to system, but some terms in common usage include:
Baseline
An approved revision of a document or source file to which subsequent changes can be made. See baselines, labels and tags.
Blame
A search for the author and revision that last modified a particular line.
Branch
A set of files under version control may be branched or forked at a point in time so that, from that time forward, two copies of those files may develop at different speeds or in different ways independently of each other.
Change
A change (or diff, or delta) represents a specific modification to a document under version control. The granularity of the modification considered a change varies between version control systems.
Change list
On many version control systems with atomic multi-change commits, a change list (or CL), change set, update, or patch identifies the set of changes made in a single commit. This can also represent a sequential view of the source code, allowing the examination of source as of any particular changelist ID.
Checkout
To check out (or co) is to create a local working copy from the repository. A user may specify a specific revision or obtain the latest. The term 'checkout' can also be used as a noun to describe the working copy. When a file has been checked out from a shared file server, it cannot be edited by other users. Think of it like a hotel, when you check out, you no longer have access to its amenities.
Clone
Cloning means creating a repository containing the revisions from another repository. This is equivalent to pushing or pulling into an empty (newly initialized) repository. As a noun, two repositories can be said to be clones if they are kept synchronized, and contain the same revisions.
Commit (noun)
Commit (verb)
To commit (check in, ci or, more rarely, install, submit or record) is to write or merge the changes made in the working copy back to the repository. A commit contains metadata, typically the author information and a commit message that describes the change.
Commit message
A short note, written by the developer, stored with the commit, which describes the commit. Ideally, it records why the modification was made, a description of the modification's effect or purpose, and non-obvious aspects of how the change works.
Conflict
A conflict occurs when different parties make changes to the same document, and the system is unable to reconcile the changes. A user must resolve the conflict by combining the changes, or by selecting one change in favour of the other.
Delta compression
Most revision control software uses delta compression, which retains only the differences between successive versions of files. This allows for more efficient storage of many different versions of files.
Dynamic stream
A stream in which some or all file versions are mirrors of the parent stream's versions.
Export
Exporting is the act of obtaining the files from the repository. It is similar to checking out except that it creates a clean directory tree without the version-control metadata used in a working copy. This is often used prior to publishing the contents, for example.
Fetch
See pull.
Forward integration
The process of merging changes made in the main trunk into a development (feature or team) branch.
Head
Also sometimes called tip, this refers to the most recent commit, either to the trunk or to a branch. The trunk and each branch have their own head, though HEAD is sometimes loosely used to refer to the trunk.
Import
Importing is the act of copying a local directory tree (that is not currently a working copy) into the repository for the first time.
Initialize
To create a new, empty repository.
Interleaved deltas
Some revision control software uses Interleaved deltas, a method that allows storing the history of text based files in a more efficient way than by using Delta compression.
Label
See tag.
Locking
When a developer locks a file, no one else can update that file until it is unlocked. Locking can be supported by the version control system, or via informal communications between developers (aka social locking).
Mainline
Similar to trunk, but there can be a mainline for each branch.
Merge
A merge or integration is an operation in which two sets of changes are applied to a file or set of files. Some sample scenarios are as follows:
A user, working on a set of files, updates or syncs their working copy with changes made, and checked into the repository, by other users.
A user tries to check in files that have been updated by others since the files were checked out, and the revision control software automatically merges the files (typically, after prompting the user if it should proceed with the automatic merge, and in some cases only doing so if the merge can be clearly and reasonably resolved).
A branch is created, the code in the files is independently edited, and the updated branch is later incorporated into a single, unified trunk.
A set of files is branched, a problem that existed before the branching is fixed in one branch, and the fix is then merged into the other branch. (This type of selective merge is sometimes known as a cherry pick to distinguish it from the complete merge in the previous case.)
Promote
The act of copying file content from a less controlled location into a more controlled location. For example, from a user's workspace into a repository, or from a stream to its parent.
Pull, push
Copy revisions from one repository into another. Pull is initiated by the receiving repository, while push is initiated by the source. Fetch is sometimes used as a synonym for pull, or to mean a pull followed by an update.
Pull request
Repository
Resolve
The act of user intervention to address a conflict between different changes to the same document.
Reverse integration
The process of merging different team branches into the main trunk of the versioning system.
Revision and version
A version is any change in form. In SVK, a Revision is the state at a point in time of the entire tree in the repository.
Share
The act of making one file or folder available in multiple branches at the same time. When a shared file is changed in one branch, it is changed in other branches.
Stream
A container for branched files that has a known relationship to other such containers. Streams form a hierarchy; each stream can inherit various properties (like versions, namespace, workflow rules, subscribers, etc.) from its parent stream.
Tag
A tag or label refers to an important snapshot in time, consistent across many files. These files at that point may all be tagged with a user-friendly, meaningful name or revision number. See baselines, labels and tags.
Trunk
The trunk is the unique line of development that is not a branch (sometimes also called Baseline, Mainline or Master)
Update
An update (or sync, but sync can also mean a combined push and pull) merges changes made in the repository (by other people, for example) into the local working copy. Update is also the term used by some CM tools (CM+, PLS, SMS) for the change package concept (see changelist). Synonymous with checkout in revision control systems that require each repository to have exactly one working copy (common in distributed systems)
Unlocking
Releasing a lock.
Working copy
The working copy is the local copy of files from a repository, at a specific time or revision. All work done to the files in a repository is initially done on a working copy, hence the name. Conceptually, it is a sandbox.
| Technology | Software development: General | null |
55973 | https://en.wikipedia.org/wiki/Outer%20ear | Outer ear | The outer ear, external ear, or auris externa is the external part of the ear, which consists of the auricle (also pinna) and the ear canal. It gathers sound energy and focuses it on the eardrum (tympanic membrane).
Structure
Auricle
The visible part is called the auricle, also known as the pinna, especially in other animals. It is composed of a thin plate of yellow elastic cartilage, covered with integument, and connected to the surrounding parts by ligaments and muscles; and to the commencement of the ear canal by fibrous tissue. Many mammals can move the pinna (with the auriculares muscles) in order to focus their hearing in a certain direction in much the same way that they can turn their eyes. Most humans do not have this ability.
Ear canal
From the pinna, the sound waves move into the ear canal (also known as the external acoustic meatus) a simple tube running through to the middle ear. This tube leads inward from the bottom of the auricula and conducts the vibrations to the tympanic cavity and amplifies frequencies in the range 2 kHz to 5 kHz.
Auricular muscles
Intrinsic muscles
The intrinsic auricular muscles are:
The helicis major is a narrow vertical band situated upon the anterior margin of the helix. It arises below, from the spina helicis, and is inserted into the anterior border of the helix, just where it is about to curve backward.
The helicis minor is an oblique fasciculus, covering the crus helicis.
The tragicus is a short, flattened vertical band on the lateral surface of the tragus. Also known as the mini lobe.
The antitragicus arises from the outer part of the antitragus, and is inserted into the cauda helicis and antihelix.
The transverse muscle is placed on the cranial surface of the pinna. It consists of scattered fibers, partly tendinous and partly muscular, extending from the eminentia conchae to the prominence corresponding with the scapha.
The oblique muscle also on the cranial surface, consists of a few fibers extending from the upper and back part of the concha to the convexity immediately above it.
The intrinsic muscles contribute to the topography of the auricle, while also function as a sphincter of the external auditory meatus. It has been suggested that during prenatal development in the womb, these muscles exert forces on the cartilage which in turn affects the shaping of the ear.
Extrinsic muscles
The extrinsic auricular muscles are the three muscles surrounding the auricula or outer ear:
anterior auricular muscle
superior auricular muscle
posterior auricular muscle
The superior muscle is the largest of the three, followed by the posterior and the anterior.
In some mammals these muscles can adjust the direction of the pinna. In humans these muscles possess very little action.
The auricularis anterior draws the auricula forward and upward, the auricularis superior slightly raises it, and the auricularis posterior draws it backward. The superior auricular muscle also acts as a stabilizer of the occipitofrontalis muscle and as a weak brow lifter. The presence of auriculomotor activity in the posterior auricular muscle causes the muscle to contract and cause the pinna to be pulled backwards and flatten when exposed to sudden, surprising sounds.
Function
One consequence of the configuration of the outer ear is selectively to boost the sound pressure 30- to 100-fold for frequencies around 3 kHz. This amplification makes humans most sensitive to frequencies in this range—and also explains why they are particularly prone to acoustical injury and hearing loss near this frequency. Most human speech sounds are also distributed in the bandwidth around 3 kHz.
Clinical significance
Malformations of the external ear can be a consequence of hereditary disease, or exposure to environmental factors such as radiation, infection. Such defects include:
A preauricular fistula, which is a long narrow tube, usually near the tragus. This can be inherited as an autosomal recessive fashion and may suffer from chronic infection in later life.
Cosmetic defects, such as very large ears, small ears.
Malformation that may lead to functional impairment, such as atresia of the external auditory meatus or aplasia of the pinna,
Genetic syndromes, which include:
Konigsmark syndrome, characterised by small ears and atresia of the external auditory canal, causing conductive hearing loss and inherited in an autosomal recessive manner.
Goldenhar syndrome, a combination of developmental abnormalities affecting the ears, eyes, bones of the skull, and vertebrae, inherited in an autosomal dominant manner.
Treacher Collins syndrome, characterised by dysplasia of the auricle, atresia of the bony part of the auditory canal, hypoplasia of the auditory ossicles and tympanic cavity, and 'mixed' deafness (both sensorineural and conductive), inherited in an autosomal dominant manner.
Crouzon syndrome, characterised by bilateral atresia of the external auditory canal, inherited in an autosomal dominant manner.
Surgery
Usually, malformations are treated with surgery, although artificial prostheses are also sometimes used.
Preauricular fistulas are generally not treated unless chronically inflamed.
Cosmetic defects without functional impairment are generally repaired after ages 6–7.
If malformations are accompanied by hearing loss amenable to correction, then the early use of hearing aids may prevent complete hearing loss.
Additional images
| Biology and health sciences | Sensory nervous system | Biology |
55982 | https://en.wikipedia.org/wiki/Dirty%20bomb | Dirty bomb | A dirty bomb or radiological dispersal device is a radiological weapon that combines radioactive material with conventional explosives. The purpose of the weapon is to contaminate the area around the dispersal agent/conventional explosion with radioactive material, serving primarily as an area denial device against civilians. It is not to be confused with a nuclear explosion, such as a fission bomb, which produces blast effects far in excess of what is achievable by the use of conventional explosives. Unlike the rain of radioactive material from a typical fission bomb, a dirty bomb's radiation can be dispersed only within a few hundred meters or a few miles of the explosion.
Dirty bombs have never been used, only tested. They are designed to disperse radioactive material over a certain area. They act through the effects of radioactive contamination on the environment and related health effects of radiation poisoning in the affected populations. The containment and decontamination of victims, as well as decontamination of the affected area require considerable time and expenses, rendering areas partly unusable and causing economic damage. Dirty bombs might be used to create mass panic as a weapon of terror.
Effect of a dirty bomb explosion
When dealing with the implications of a dirty bomb attack, there are two main areas to be addressed: the civilian impact, not only dealing with immediate casualties and long term health issues, but also the psychological effect, and the economic impact. With no prior event of a dirty bomb detonation, it is considered difficult to predict the impact. Several analyses have predicted that radiological dispersal devices will neither sicken nor kill many people.
Differences between dirty bombs and fission bombs
Source: Adapted from Levi MA, Kelly HC. "Weapons of mass disruption". Sci Am. 2002 Nov;287(5):76-81.
Accidents with radioactives
The effects of uncontrolled radioactive contamination have been reported several times.
One example is the radiological accident occurring in Goiânia, Brazil, between September 1987 and March 1988: Two metal scavengers broke into an abandoned radiotherapy clinic and removed a teletherapy source capsule containing powdered caesium-137 with an activity of 50 TBq. They brought it back to the home of one of the men to take it apart and sell as scrap metal. Later that day both men were showing acute signs of radiation illness with vomiting and one of them had a swollen hand and diarrhea. A few days later one of the men punctured the thick window of the capsule, allowing the caesium chloride powder to leak out and when realizing the powder glowed blue in the dark, brought it back home to his family and friends to show it off. After two weeks of spread by contact contamination causing an increasing number of adverse health effects, the correct diagnosis of acute radiation sickness was made at a hospital and proper precautions could be put into procedure. By this time 249 people were contaminated, 151 exhibited both external and internal contamination, of whom 20 people were seriously ill and five people died.
The Goiânia incident to some extent predicts the contamination pattern if it is not immediately realized that the explosion spread radioactive material, but also how fatal even very small amounts of ingested radioactive powder can be. This raises worries of terrorists using powdered alpha emitting material, that if ingested can pose a serious health risk, as in the case of Alexander Litvinenko, who was poisoned by tea with polonium-210. "Smoky bombs" based on alpha emitters might be just as dangerous as beta or gamma emitting dirty bombs.
Public perception of risks
Although the exposure might be minimal, many people find radiation exposure especially frightening because it is something they cannot see or feel, and it therefore becomes an unknown source of danger. When United States Attorney General John Ashcroft on June 10, 2002, announced the arrest of José Padilla, allegedly plotting to detonate such a weapon, he said:
This public fear of radiation also plays a big role in why the costs of a radiological dispersal device impact on a major metropolitan area (such as lower Manhattan) might be equal to or even larger than that of the 9/11 attacks. Assuming the radiation levels are not too high and the area does not need to be abandoned such as the town of Pripyat near the Chernobyl reactor, an expensive and time-consuming cleanup procedure will begin. This will mainly consist of tearing down highly contaminated buildings, digging up contaminated soil and quickly applying sticky substances to remaining surfaces so that radioactive particles adhere before radioactivity penetrates the building materials. These procedures are the current state of the art for radioactive contamination cleanup, but some experts say that a complete cleanup of external surfaces in an urban area to current decontamination limits may not be technically feasible. Loss of working hours will be vast during cleanup, but even after the radiation levels reduce to an acceptable level, there might be residual public fear of the site including possible unwillingness to conduct business as usual in the area. Tourist traffic is likely never to resume.
Dirty bombs and terrorism
Since the 9/11 attacks, the fear of terrorist groups using dirty bombs has increased, which has been frequently reported in the media. The meaning of terrorism used here, is described by the U.S. Department of Defense's definition, which is "the calculated use of unlawful violence or threat of unlawful violence to inculcate fear; intended to coerce or to intimidate governments or societies in the pursuit of goals that are generally political, religious, or ideological."
Constructing and obtaining material for a dirty bomb
In order for a terrorist organization to construct and detonate a dirty bomb, it must acquire radioactive material. Possible radiological dispersal device material could come from the millions of radioactive sources used worldwide in the industry, for medical purposes and in academic applications mainly for research. Of these sources, only nine reactor-produced isotopes stand out as being suitable for radiological terror: americium-241, californium-252, caesium-137, cobalt-60, iridium-192, plutonium-238, polonium-210, radium-226 and strontium-90, and even from these it is possible that radium-226 and polonium-210 do not pose a significant threat. Of these sources the U.S. Nuclear Regulatory Commission has estimated that within the U.S., approximately one source is lost, abandoned or stolen every day of the year. Within the European Union the annual estimate is 70. There exist thousands of such "orphan" sources scattered throughout the world, but of those reported lost, no more than an estimated 20 percent can be classified as potential high security concerns if used in a radiological dispersal device. Russia is believed to house thousands of orphan sources, which were lost following the collapse of the Soviet Union. A large but unknown number of these sources probably belong to the high security risk category. These include the beta-emitting strontium-90 sources used as radioisotope thermoelectric generators for beacons in lighthouses in remote areas of Russia. In December 2001, three Georgian woodcutters stumbled over such a power generator and dragged it back to their camp site to use it as a heat source. Within hours they suffered from acute radiation sickness and sought hospital treatment. The International Atomic Energy Agency (IAEA) later stated that it contained approximately of strontium, equivalent to the amount of radioactivity released immediately after the Chernobyl accident (though the total radioactivity release from Chernobyl was 2500 times greater at around ).
Although a terrorist organization might obtain radioactive material through the "black market", and there has been a steady increase in illicit trafficking of radioactive sources from 1996 to 2004, these recorded trafficking incidents mainly refer to rediscovered orphan sources without any sign of criminal activity, and it has been argued that there is no conclusive evidence for such a market. In addition to the hurdles of obtaining usable radioactive material, there are several conflicting requirements regarding the properties of the material the terrorists need to take into consideration: First, the source should be "sufficiently" radioactive to create direct radiological damage at the explosion or at least to perform societal damage or disruption. Second, the source should be transportable with enough shielding to protect the carrier, but not so much that it will be too heavy to maneuver. Third, the source should be sufficiently dispersible to effectively contaminate the area around the explosion.
Possibility of use by terrorist groups
The first attempt of radiological terror was reportedly carried out in November 1995 by a group of Chechen separatists, who buried a caesium-137 source wrapped in explosives at the Izmaylovsky Park in Moscow. A Chechen rebel leader alerted the media, the bomb was never activated, and the incident amounted to a mere publicity stunt. In December 1998, a second attempt was announced by the Chechen Security Service, who discovered a container filled with radioactive materials attached to an explosive mine. The bomb was hidden near a railway line in the suburban area Argun, ten miles east of the Chechen capital of Grozny. The same Chechen separatist group was suspected to be involved.
On 8 May 2002, José Padilla (a.k.a. Abdulla al-Muhajir) was arrested on suspicion that he was an al-Qaeda terrorist planning to detonate a dirty bomb in the U.S. This suspicion was raised by information obtained from an arrested terrorist in U.S. custody, Abu Zubaydah, who under interrogation revealed that the organization was close to constructing a dirty bomb. Although Padilla had not obtained radioactive material or explosives at the time of arrest, law enforcement authorities uncovered evidence that he was on reconnaissance for usable radioactive material and possible locations for detonation. It has been doubted whether José Padilla was preparing such an attack, and it has been claimed that the arrest was highly politically motivated, given the pre-9/11 security lapses by the CIA and FBI.
In 2006, Dhiren Barot from North London pleaded guilty of conspiring to murder people in the United Kingdom and United States using a radioactive dirty bomb. He planned to target underground car parks within the UK and buildings in the U.S. such as the International Monetary Fund, World Bank buildings in Washington D.C., the New York Stock Exchange, Citigroup buildings and the Prudential Financial buildings in Newark, New Jersey. He also faces 12 other charges including, conspiracy to commit public nuisance, seven charges of making a record of information for terrorist purposes and four charges of possessing a record of information for terrorist purposes. Experts say if the plot to use the dirty bomb was carried out "it would have been unlikely to cause deaths, but was designed to affect about 500 people".
In January 2009, a leaked FBI report described the results of a search of the Maine home of James G. Cummings, a white supremacist who had been shot and killed by his wife. Investigators found four one-gallon containers of 35 percent hydrogen peroxide, uranium, thorium, lithium metal, aluminium powder, beryllium, boron, black iron oxide and magnesium as well as literature on how to build dirty bombs and information about caesium-137, strontium-90 and cobalt-60, radioactive materials. Officials confirmed the veracity of the report but stated that the public was never at risk.
In July 2014, ISIS militants seized of uranium compounds from Mosul University. The material was unenriched and so could not be used to build a conventional fission bomb, but a dirty bomb is a theoretical possibility. Nonetheless, uranium's relatively low radioactivity makes it a poor candidate for use in a dirty bomb.
Terrorist organizations may also capitalize on the fear of radiation to create weapons of mass disruption rather than weapons of mass destruction. A fearful public response may in itself accomplish the goals of a terrorist organization to gain publicity or destabilize society. Even simply stealing radioactive materials may trigger a panic reaction from the general public. Similarly, a small-scale release of radioactive materials or a threat of such a release may be considered sufficient for a terror attack. Particular concern is directed towards the medical sector and healthcare sites, which are "intrinsically more vulnerable than conventional licensed nuclear sites". Opportunistic attacks may range to even kidnapping patients whose treatment involve radioactive materials. In the Goiânia accident, over 100,000 people admitted themselves to monitoring, while only 49 were admitted to hospitals. Other benefits to a terrorist organization of a dirty bomb include economic disruption in the area affected, abandonment of affected assets (such a buildings, subways) due to public concern, and international publicity useful for recruitment.
Tests
Israel carried out a four-year series of tests on nuclear explosives to measure the effects were hostile forces ever to use them against Israel, Haaretz reported in 2015. According to the report, high-level radiation was measured only at the center of the explosions, while the level of dispersal of radiation by particles carried by the wind (fallout) was low. The bombs reportedly did not pose a significant danger beyond their psychological effect.
Detection and prevention
Dirty bombs may be prevented by detecting illicit radioactive materials in shipping with tools such as a Radiation Portal Monitor. Similarly, unshielded radioactive materials may be detected at checkpoints by Geiger counters, gamma-ray detectors, and even Customs and Border Patrol (CBS) pager-sized radiation detectors. Hidden materials may also be detected by x-ray inspection and heat emitted may be picked up by infrared detectors. Such devices, however, may be circumvented by simply transporting materials across unguarded stretches of coastline or other barren border areas.
One proposed method for detecting shielded Dirty Bombs is Nanosecond Neutron Analysis (NNA). Designed originally for the detection of explosives and hazardous chemicals, NNA is also applicable to fissile materials. NNA determines what chemicals are present in an investigated device by analyzing emitted γ-emission neutrons and α-particles created from a reaction in the neutron generator. The system records the temporal and spatial displacement of the neutrons and α-particles within separate 3D regions. A prototype dirty-bomb detection device created with NNA is demonstrated to be able to detect uranium from behind a 5 cm-thick lead wall. Other radioactive material detectors include Radiation Assessment and Identification (RAID) and Sensor for Measurement and Analysis of Radiation Transients, both developed by Sandia National Laboratories. Sodium iodide scintillator based aerial radiation detection systems are capable to detect International Atomic Energy Agency (IAEA) defined dangerous quantities of radioactive material and have been deployed by the New York City Police Department (NYPD) Counterterrorism Bureau.
The IAEA recommends certain devices be used in tandem at country borders to prevent transfer of radioactive materials, and thus the building of dirty bombs. They define the four main goals of radiation detection instruments as detection, verification, assessment and localization, and identification as a means to escalate a potential radiological situation. The IAEA also defines the following types of instruments:
Pocket-Type Instruments: these instruments provide a low-power, mobile option to detection that allows for security officers to passively scan an area for radioactive materials. These devices should be easily worn, should have an alarm threshold of three times normal radiation levels, and should have a long battery life - over 800 hours.
Handheld Instruments: these instruments may be used to detect all types of radiation (including neutron) and may be used to search specific targets flexibly. These instruments should aim for ease of use and speed, ideally weighing less than 2 kg and being able to make measurements in less than a second.
Fixed, installed instruments: these instruments provide a continuous, automatic detection system that can monitor pedestrians and vehicles that pass through. To work effectively pedestrians and vehicles should be led close to the detectors, as performance is directly related to range.
Legislative and regulatory actions can also be used to prevent access to materials needed to create a dirty bomb. Examples include the 2006 U.S. Dirty Bomb Bill, the Yucca Flats proposal, and the Nunn-Lungar act. Similarly, close monitoring and restrictions of radioactive materials may provide security for materials in vulnerable private-sector applications, most notably in the medical sector where such materials are used for treatments. Suggestions for increased security include isolation of materials in remote locations and strict limitation of access.
One way to mitigate a major effect of a radiological weapons may also be to educate the public on the nature of radioactive materials. As one of the major concerns of a dirty bomb is the public panic proper education may prove a viable counter-measure. Education on radiation is considered by some to be "the most neglected issue related to radiological terrorism".
Personal safety
The dangers of a dirty bomb come from the initial blast and the radioactive materials To mitigate the risk of radiation exposure, FEMA suggests the following guidelines:
Covering the mouth/nose with cloth to reduce risk of breathing in radioactive materials.
Avoiding touching materials touched by the explosion.
Quickly relocating inside to shield from radiation.
Remove and pack up clothes. Keep clothes until instructed by authorities how to dispose of them.
Keep radioactive dust outside.
Remove all dust possible by showering with soap and water.
Avoid taking potassium iodide, as it only prevents effects from radioactive iodine and may instead cause a dangerous reaction.
Treatment
, research is under way to find radioactive decontanimation drugs to remove radioactive elements from the body. One drug candidate under investigation is HOPO 14-1.
Other uses of the term
The term has also been used historically to refer to certain types of nuclear weapons. Due to the inefficiency of early nuclear weapons, only a small amount of the nuclear material would be consumed during the explosion. Little Boy had an efficiency of only 1.4%. Fat Man, which used a different design and a different fissile material, had an efficiency of 14%. Thus, they tended to disperse large amounts of unused fissile material, and the fission products, which are on average much more dangerous, in the form of nuclear fallout. During the 1950s, there was considerable debate over whether "clean" bombs could be produced and these were often contrasted with "dirty" bombs. "Clean" bombs were often a stated goal and scientists and administrators said that high-efficiency nuclear weapon design could create explosions that generated almost all of their energy in the form of nuclear fusion, which does not create harmful fission products.
But the Castle Bravo accident of 1954, in which a thermonuclear weapon produced a large amount of fallout that was dispersed among human populations, suggested that this was not what was actually being used in modern thermonuclear weapons, which derive around half of their yield from a final fission stage of the fast fissioning of the uranium tamper of the secondary. While some proposed producing "clean" weapons, other theorists noted that one could make a nuclear weapon intentionally "dirty" by "salting" it with a material, which would generate large amounts of long-lasting fallout when irradiated by the weapon core. These are known as salted bombs; a specific subtype often noted is a cobalt bomb.
In popular culture
In the 2018 video game Detroit: Become Human, numerous endings depict Markus, one of three playable android characters in the game, setting off a cobalt-derived dirty bomb in southern Detroit to force the retreat of authorities.
| Technology | Weapon of mass destruction | null |
55983 | https://en.wikipedia.org/wiki/Black%20rat | Black rat | The black rat (Rattus rattus), also known as the roof rat, ship rat, or house rat, is a common long-tailed rodent of the stereotypical rat genus Rattus, in the subfamily Murinae. It likely originated in the Indian subcontinent, but is now found worldwide.
The black rat is black to light brown in colour with a lighter underside. It is a generalist omnivore and a serious pest to farmers because it feeds on a wide range of agricultural crops. It is sometimes kept as a pet. In parts of India, it is considered sacred and respected in the Karni Mata Temple in Deshnoke.
Taxonomy
Mus rattus was the scientific name proposed by Carl Linnaeus in 1758 for the black rat.
Three subspecies were once recognized, but today are considered invalid and are now known to be actually color morphs:
Rattus rattus rattus – roof rat
Rattus rattus alexandrinus – Alexandrine rat
Rattus rattus frugivorus – fruit rat
Characteristics
A typical adult black rat is long, not including a tail, and weighs , depending on the subspecies. Black rats typically live for about one year in the wild and up to four years in captivity. Despite its name, the black rat exhibits several colour forms. It is usually black to light brown in colour with a lighter underside. In England during the 1920s, several variations were bred and shown alongside domesticated brown rats. This included an unusual green-tinted variety.
Origin
The black rat was present in prehistoric Europe and in the Levant during postglacial periods. The black rat in the Mediterranean region differs genetically from its South Asian ancestor by having 38 instead of 42 chromosomes. Its closest relative is the Asian house rat (R. tanezumi) from Southeast Asia. The two diverged about 120,000 years ago in southwestern Asia. It is unclear how the rat made its way to Europe due to insufficient data, although a land route seems more likely based on the distribution of European haplogroup "A". The black rat spread throughout Europe with the Roman conquest, but declined around the 6th century, possibly due to collapse of the Roman grain trade, climate cooling, or the Justinianic Plague. A genetically different rat population of haplogroup A replaced the Roman population in the medieval times in Europe.
It is a resilient vector for many diseases because of its ability to hold so many infectious bacteria in its blood. It was formerly thought to have played a primary role in spreading bacteria contained in fleas on its body, such as the plague bacterium (Yersinia pestis) which is responsible for the Plague of Justinian and the Black Death. However, recent studies have called this theory into question and instead posit humans themselves as the vector, as the movements of the epidemics and the black rat populations do not show historical or geographical correspondence. A study published in 2015 indicates that other Asiatic rodents served as plague reservoirs, from which infections spread as far west as Europe via trade routes, both overland and maritime. Although the black rat was certainly a plague vector in European ports, the spread of the plague beyond areas colonized by rats suggests that the plague was also circulated by humans after reaching Europe.
Distribution and habitat
The black rat originated in India and Southeast Asia, and spread to the Near East and Egypt, and then throughout the Roman Empire, reaching Great Britain as early as the 1st century AD. Europeans subsequently spread it throughout the world. The black rat is again largely confined to warmer areas, having been supplanted by the brown rat (Rattus norvegicus) in cooler regions and urban areas. In addition to the brown rat being larger and more aggressive, the change from wooden structures and thatched roofs to bricked and tiled buildings favored the burrowing brown rats over the arboreal black rats. In addition, brown rats eat a wider variety of foods, and are more resistant to weather extremes.
Black rat populations can increase exponentially under certain circumstances, perhaps having to do with the timing of the fruiting of the bamboo plant, and cause devastation to the plantings of subsistence farmers; this phenomenon is known as mautam in parts of India.
Black rats are thought to have arrived in Australia with the First Fleet, and subsequently spread to many coastal regions in the country.
Black rats adapt to a wide range of habitats. In urban areas they are found around warehouses, residential buildings, and other human settlements. They are also found in agricultural areas, such as in barns and crop fields. In urban areas, they prefer to live in dry upper levels of buildings, so they are commonly found in wall cavities and false ceilings. In the wild, black rats live in cliffs, rocks, the ground, and trees. They are great climbers and prefer to live in palms and trees, such as pine trees. Their nests are typically spherical and made of shredded material, including sticks, leaves, other vegetation and cloth. In the absence of palms or trees, they can burrow into the ground. Black rats are also found around fences, ponds, riverbanks, streams, and reservoirs.
Behaviour and ecology
It is thought that male and female rats have similarly sized home ranges during the winter, but male rats increase the size of their home range during the breeding season. Along with differing between rats of different sex, home range also differs depending on the type of forest in which the black rat inhabits. For example, home ranges in the southern beech forests of the South Island, New Zealand appear to be much larger than the non-beech forests of the North Island. Due to the limited number of rats that are studied in home range studies, the estimated sizes of rat home ranges in different rat demographic groups are inconclusive.
Diet and foraging
Black rats are considered omnivores and eat a wide range of foods, including seeds, fruit, stems, leaves, fungi, and a variety of invertebrates and vertebrates. They are generalists, and thus not very specific in their food preferences, which is indicated by their tendency to feed on any meal provided for cows, swine, chickens, cats and dogs. They are similar to the tree squirrel in their preference of fruits and nuts. They eat about per day and drink about per day. Their diet is high in water content. They are a threat to many natural habitats because they feed on birds and insects. They are also a threat to many farmers, since they feed on a variety of agricultural-based crops, such as cereals, sugar cane, coconuts, cocoa, oranges, and coffee beans.
The black rat displays flexibility in its foraging behaviour. It is a predatory species and adapts to different micro-habitats. It often meets and forages together in close proximity within and between sexes. It tends to forage after sunset. If the food cannot be eaten quickly, it searches for a place to carry and hoard to eat at a later time. Although it eats a broad range of foods, it is a highly selective feeder; only a restricted selection of the foods is dominating. When offered a wide diversity of foods, it eats only a small sample of each. This allows it to monitor the quality of foods that are present year round, such as leaves, as well as seasonal foods, such as herbs and insects. This method of operating on a set of foraging standards ultimately determines the final composition of its meals. Also, by sampling the available food in an area, it maintains a dynamic food supply, balance its nutrient intake, and avoids intoxication by secondary compounds.
Nesting behaviour
Through the usage of tracking devices such as radio transmitters, rats have been found to occupy dens located in trees, as well as on the ground. In Puketi Forest in the Northland Region of New Zealand, rats have been found to form dens together. Rats appear to den and forage in separate areas in their home range depending on the availability of food resources. Research shows that, in New South Wales, the black rat prefers to inhabit lower leaf litter of forest habitat. There is also an apparent correlation between the canopy height and logs and the presence of black rats. This correlation may be a result of the distribution of the abundance of prey as well as available refuges for rats to avoid predators. As found in North Head, New South Wales, there is positive correlation between rat abundance, leaf litter cover, canopy height, and litter depth. All other habitat variables showed little to no correlation. While this species' relative, the brown (Norway) rat, prefers to nest near the ground of a building the black rat will prefer the upper floors and roof. Because of this habit they have been given the common name roof rat.
Diseases
Black rats (or their ectoparasites) can carry a number of pathogens, of which bubonic plague (via the Oriental rat flea), typhus, Weil's disease, toxoplasmosis and trichinosis are the best known. It has been hypothesized that the displacement of black rats by brown rats led to the decline of the Black Death. This theory has, however, been deprecated, as the dates of these displacements do not match the increases and decreases in plague outbreaks.
Rats serve as outstanding vectors for transmittance of diseases because they can carry bacteria and viruses in their systems. A number of bacterial diseases are common to rats, and these include Streptococcus pneumoniae, Corynebacterium kutsheri, Bacillus piliformis, Pasteurella pneumotropica, and Streptobacillus moniliformis, to name a few. All of these bacteria are disease causing agents in humans. In some cases, these diseases are incurable.
Predators
The black rat is prey to cats and owls in domestic settings. In less urban settings, rats are preyed on by weasels, foxes and coyotes. These predators have little effect on the control of the black rat population because black rats are agile and fast climbers. In addition to agility, the black rat also uses its keen sense of hearing to detect danger and quickly evade mammalian and avian predators.
As an invasive species
Damage caused
After Rattus rattus was introduced into the northern islands of New Zealand, they fed on the seedlings, adversely affecting the ecology of the islands. Even after eradication of R. rattus, the negative effects may take decades to reverse. When consuming these seabirds and seabird eggs, these rats reduce the pH of the soil. This harms plant species by reducing nutrient availability in soil, thus decreasing the probability of seed germination. For example, research conducted by Hoffman et al. indicates a large impact on 16 indigenous plant species directly preyed on by R. rattus. These plants displayed a negative correlation in germination and growth in the presence of black rats.
Rats prefer to forage in forest habitats. In the Ogasawara islands, they prey on the indigenous snails and seedlings. Snails that inhabit the leaf litter of these islands showed a significant decline in population on the introduction of Rattus rattus. The black rat shows a preference for snails with larger shells (greater than 10 mm), and this led to a great decline in the population of snails with larger shells. A lack of prey refuges makes it more difficult for the snail to avoid the rat.
Complex pest
The black rat is a complex pest, defined as one that influences the environment in both harmful and beneficial ways. In many cases, after the black rat is introduced into a new area, the population size of some native species declines or goes extinct. This is because the black rat is a good generalist with a wide dietary niche and a preference for complex habitats; this causes strong competition for resources among small animals. This has led to the black rat completely displacing many native species in Madagascar, the Galapagos, and the Florida Keys. In a study by Stokes et al., habitats suitable for the native bush rat, Rattus fuscipes, of Australia are often invaded by the black rat and are eventually occupied by only the black rat. When the abundances of these two rat species were compared in different micro-habitats, both were found to be affected by micro-habitat disturbances, but the black rat was most abundant in areas of high disturbance; this indicates it has a better dispersal ability.
Despite the black rat's tendency to displace native species, it can also aid in increasing species population numbers and maintaining species diversity. The bush rat, a common vector for spore dispersal of truffles, has been extirpated from many micro-habitats of Australia. In the absence of a vector, the diversity of truffle species would be expected to decline. In a study in New South Wales, Australia it was found that, although the bush rat consumes a diversity of truffle species, the black rat consumes as much of the diverse fungi as the natives and is an effective vector for spore dispersal. Since the black rat now occupies many of the micro-habitats that were previously inhabited by the bush rat, the black rat plays an important ecological role in the dispersal of fungal spores. By eradicating the black rat populations in Australia, the diversity of fungi would decline, potentially doing more harm than good.
Control methods
Large-scale rat control programs have been taken to maintain a steady level of the invasive predators in order to conserve the native species in New Zealand such as kokako and mohua. Pesticides, such as pindone and 1080 (sodium fluoroacetate), are commonly distributed via aerial spray by helicopter as a method of mass control on islands infested with invasive rat populations. Bait, such as brodifacoum, is also used along with coloured dyes (used to deter birds from eating the baits) in order to kill and identify rats for experimental and tracking purposes. Another method to track rats is the use of wired cage traps, which are used along with bait, such as rolled oats and peanut butter, to tag and track rats to determine population sizes through methods like mark-recapture and radio-tracking. Tracking tunnels (coreflute tunnels containing an inked card) are also commonly used monitoring devices, as are chew-cards containing peanut butter. Poison control methods are effective in reducing rat populations to nonthreatening sizes, but rat populations often rebound to normal size within months. Besides their highly adaptive foraging behavior and fast reproduction, the exact mechanisms for their rebound is unclear and are still being studied.
In 2010, the Sociedad Ornitológica Puertorriqueña (Puerto Rican Bird Society) and the Ponce Yacht and Fishing Club launched a campaign to eradicate the black rat from the Isla Ratones (Mice Island) and Isla Cardona (Cardona Island) islands off the municipality of Ponce, Puerto Rico.
Decline in population
Eradication projects have eliminated black rats from Lundy in the Bristol Channel (2006) and from the Shiant Islands in the Outer Hebrides (2016). Populations probably survive on other islands (e.g. Inchcolm) and in localised areas of the British mainland. Recent National Biodiversity Network data show populations around the U.K., particularly in ports and port towns.
| Biology and health sciences | Rodents | Animals |
55999 | https://en.wikipedia.org/wiki/Tongue | Tongue | The tongue is a muscular organ in the mouth of a typical tetrapod. It manipulates food for chewing and swallowing as part of the digestive process, and is the primary organ of taste. The tongue's upper surface (dorsum) is covered by taste buds housed in numerous lingual papillae. It is sensitive and kept moist by saliva and is richly supplied with nerves and blood vessels. The tongue also serves as a natural means of cleaning the teeth. A major function of the tongue is to enable speech in humans and vocalization in other animals.
The human tongue is divided into two parts, an oral part at the front and a pharyngeal part at the back. The left and right sides are also separated along most of its length by a vertical section of fibrous tissue (the lingual septum) that results in a groove, the median sulcus, on the tongue's surface.
There are two groups of glossal muscles. The four intrinsic muscles alter the shape of the tongue and are not attached to bone. The four paired extrinsic muscles change the position of the tongue and are anchored to bone.
Etymology
The word tongue derives from the Old English tunge, which comes from Proto-Germanic *tungōn. It has cognates in other Germanic languages—for example tonge in West Frisian, tong in Dutch and Afrikaans, Zunge in German, tunge in Danish and Norwegian, and tunga in Icelandic, Faroese and Swedish. The ue ending of the word seems to be a fourteenth-century attempt to show "proper pronunciation", but it is "neither etymological nor phonetic". Some used the spelling tunge and tonge as late as the sixteenth century.
In humans
Structure
The tongue is a muscular hydrostat that forms part of the floor of the oral cavity. The left and right sides of the tongue are separated by a vertical section of fibrous tissue known as the lingual septum. This division is along the length of the tongue save for the very back of the pharyngeal part and is visible as a groove called the median sulcus. The human tongue is divided into anterior and posterior parts by the terminal sulcus, which is a "V"-shaped groove. The apex of the terminal sulcus is marked by a blind foramen, the foramen cecum, which is a remnant of the median thyroid diverticulum in early embryonic development. The anterior oral part is the visible part situated at the front and makes up roughly two-thirds the length of the tongue. The posterior pharyngeal part is the part closest to the throat, roughly one-third of its length. These parts differ in terms of their embryological development and nerve supply.
The anterior tongue is, at its apex, thin and narrow. It is directed forward against the lingual surfaces of the lower incisor teeth. The posterior part is, at its root, directed backward, and connected with the hyoid bone by the hyoglossi and genioglossi muscles and the hyoglossal membrane, with the epiglottis by three glossoepiglottic folds of mucous membrane, with the soft palate by the glossopalatine arches, and with the pharynx by the superior pharyngeal constrictor muscle and the mucous membrane. It also forms the anterior wall of the oropharynx.
The average length of the human tongue from the oropharynx to the tip is 10 cm. The average weight of the human tongue from adult males is 99g and for adult females 79g.
In phonetics and phonology, a distinction is made between the tip of the tongue and the blade (the portion just behind the tip). Sounds made with the tongue tip are said to be apical, while those made with the tongue blade are said to be laminal.
Upper surface
The upper surface of the tongue is called the dorsum, and is divided by a groove into symmetrical halves by the median sulcus. The foramen cecum marks the end of this division (at about 2.5 cm from the root of the tongue) and the beginning of the terminal sulcus. The foramen cecum is also the point of attachment of the thyroglossal duct and is formed during the descent of the thyroid diverticulum in embryonic development.
The terminal sulcus is a shallow groove that runs forward as a shallow groove in a V shape from the foramen cecum, forwards and outwards to the margins (borders) of the tongue. The terminal sulcus divides the tongue into a posterior pharyngeal part and an anterior oral part. The pharyngeal part is supplied by the glossopharyngeal nerve and the oral part is supplied by the lingual nerve (a branch of the mandibular branch (V3) of the trigeminal nerve) for somatosensory perception and by the chorda tympani (a branch of the facial nerve) for taste perception.
Both parts of the tongue develop from different pharyngeal arches.
Undersurface
On the undersurface of the tongue is a fold of mucous membrane called the frenulum that tethers the tongue at the midline to the floor of the mouth. On either side of the frenulum are small prominences called sublingual caruncles that the major salivary submandibular glands drain into.
Muscles
The eight muscles of the human tongue are classified as either intrinsic or extrinsic. The four intrinsic muscles act to change the shape of the tongue, and are not attached to any bone. The four extrinsic muscles act to change the position of the tongue, and are anchored to bone.
Extrinsic
The four extrinsic muscles originate from bone and extend to the tongue. They are the genioglossus, the hyoglossus (often including the chondroglossus) the styloglossus, and the palatoglossus. Their main functions are altering the tongue's position allowing for protrusion, retraction, and side-to-side movement.
The genioglossus arises from the mandible and protrudes the tongue. It is also known as the tongue's "safety muscle" since it is the only muscle that propels the tongue forward.
The hyoglossus, arises from the hyoid bone and retracts and depresses the tongue. The chondroglossus is often included with this muscle.
The styloglossus arises from the styloid process of the temporal bone and draws the sides of the tongue up to create a trough for swallowing.
The palatoglossus arises from the palatine aponeurosis, and depresses the soft palate, moves the palatoglossal fold towards the midline, and elevates the back of the tongue during swallowing.
Intrinsic
Four paired intrinsic muscles of the tongue originate and insert within the tongue, running along its length. They are the superior longitudinal muscle, the inferior longitudinal muscle, the vertical muscle, and the transverse muscle. These muscles alter the shape of the tongue by lengthening and shortening it, curling and uncurling its apex and edges as in tongue rolling, and flattening and rounding its surface. This provides shape and helps facilitate speech, swallowing, and eating.
The superior longitudinal muscle runs along the upper surface of the tongue under the mucous membrane, and functions to shorten and curl the tongue upward. It originates near the epiglottis, at the hyoid bone, from the median fibrous septum.
The inferior longitudinal muscle lines the sides of the tongue, and is joined to the styloglossus muscle. It functions to shorten and curl the tongue downward.
The vertical muscle is located in the middle of the tongue, and joins the superior and inferior longitudinal muscles. It functions to flatten the tongue.
The transverse muscle divides the tongue at the middle, and is attached to the mucous membranes that run along the sides. It functions to lengthen and narrow the tongue.
Blood supply
The tongue receives its blood supply primarily from the lingual artery, a branch of the external carotid artery. The lingual veins drain into the internal jugular vein. The floor of the mouth also receives its blood supply from the lingual artery. There is also a secondary blood supply to the root of tongue from the tonsillar branch of the facial artery and the ascending pharyngeal artery.
An area in the neck sometimes called the Pirogov triangle is formed by the intermediate tendon of the digastric muscle, the posterior border of the mylohyoid muscle, and the hypoglossal nerve. The lingual artery is a good place to stop severe hemorrhage from the tongue.
Nerve supply
Innervation of the tongue consists of motor fibers, special sensory fibers for taste, and general sensory fibers for sensation.
Motor supply for all intrinsic and extrinsic muscles of the tongue is supplied by efferent motor nerve fibers from the hypoglossal nerve (CN XII), with the exception of the palatoglossus, which is innervated by the vagus nerve (CN X).
Innervation of taste and sensation is different for the anterior and posterior part of the tongue because they are derived from different embryological structures (pharyngeal arch 1 and pharyngeal arches 3 and 4, respectively).
Anterior two-thirds of tongue (anterior to the vallate papillae):
Taste: chorda tympani branch of the facial nerve (CN VII) via special visceral afferent fibers
Sensation: lingual branch of the mandibular (V3) division of the trigeminal nerve (CN V) via general visceral afferent fibers
Posterior one third of tongue:
Taste and sensation: glossopharyngeal nerve (CN IX) via a mixture of special and general visceral afferent fibers
Base of tongue
Taste and sensation: internal branch of the superior laryngeal nerve (itself a branch of the vagus nerve, CN X)
Lymphatic drainage
The tip of tongue drains to the submental nodes. The left and right halves of the anterior two-thirds of the tongue drains to submandibular lymph nodes, while the posterior one-third of the tongue drains to the jugulo-omohyoid nodes.
Microanatomy
The upper surface of the tongue is covered in masticatory mucosa, a type of oral mucosa, which is of keratinized stratified squamous epithelium. Embedded in this are numerous papillae, some of which house the taste buds and their taste receptors. The lingual papillae consist of filiform, fungiform, vallate and foliate papillae, and only the filiform papillae are not associated with any taste buds.
The tongue can divide itself in dorsal and ventral surface. The dorsal surface is a stratified squamous keratinized epithelium, which is characterized by numerous mucosal projections called papillae. The lingual papillae covers the dorsal side of the tongue towards the front of the terminal groove. The ventral surface is stratified squamous non-keratinized epithelium which is smooth.
Development
The tongue begins to develop in the fourth week of embryonic development from a median swelling – the median tongue bud (tuberculum impar) of the first pharyngeal arch.
In the fifth week a pair of lateral lingual swellings, one on the right side and one on the left, form on the first pharyngeal arch. These lingual swellings quickly expand and cover the median tongue bud. They form the anterior part of the tongue that makes up two-thirds of the length of the tongue, and continue to develop through prenatal development. The line of their fusion is marked by the median sulcus.
In the fourth week, a swelling appears from the second pharyngeal arch, in the midline, called the copula. During the fifth and sixth weeks, the copula is overgrown by a swelling from the third and fourth arches (mainly from the third arch) called the hypopharyngeal eminence, and this develops into the posterior part of the tongue (the other third and the posterior most part of the tongue is developed from the fourth pharyngeal arch). The hypopharyngeal eminence develops mainly by the growth of endoderm from the third pharyngeal arch. The boundary between the two parts of the tongue, the anterior from the first arch and the posterior from the third arch is marked by the terminal sulcus. The terminal sulcus is shaped like a V with the tip of the V situated posteriorly. At the tip of the terminal sulcus is the foramen cecum, which is the point of attachment of the thyroglossal duct where the embryonic thyroid begins to descend.
Function
Taste
Chemicals that stimulate taste receptor cells are known as tastants. Once a tastant is dissolved in saliva, it can make contact with the plasma membrane of the gustatory hairs, which are the sites of taste transduction.
The tongue is equipped with many taste buds on its dorsal surface, and each taste bud is equipped with taste receptor cells that can sense particular classes of tastes. Distinct types of taste receptor cells respectively detect substances that are sweet, bitter, salty, sour, spicy, or taste of umami. Umami receptor cells are the least understood and accordingly are the type most intensively under research. There is a common misconception that different sections of the tongue are exclusively responsible for different basic tastes. Although widely taught in schools in the form of the tongue map, this is incorrect; all taste sensations come from all regions of the tongue, although certain parts are more sensitive to certain tastes.
Mastication
The tongue is an important accessory organ in the digestive system. The tongue is used for crushing food against the hard palate, during mastication and manipulation of food for softening prior to swallowing. The epithelium on the tongue's upper, or dorsal surface is keratinised. Consequently, the tongue can grind against the hard palate without being itself damaged or irritated.
Speech
The tongue is one of the primary articulators in the production of speech, and this is facilitated by both the extrinsic muscles that move the tongue and the intrinsic muscles that change its shape. Specifically, different vowels are articulated by changing the tongue's height and retraction to alter the resonant properties of the vocal tract. These resonant properties amplify specific harmonic frequencies (formants) that are different for each vowel, while attenuating other harmonics. For example, [a] is produced with the tongue lowered and centered and [i] is produced with the tongue raised and fronted. Consonants are articulated by constricting airflow through the vocal tract, and many consonants feature a constriction between the tongue and some other part of the vocal tract. For example, alveolar consonants like [s] and [n] are articulated with the tongue against the alveolar ridge, while velar consonants like [k] and [g] are articulated with the tongue dorsum against the soft palate (velum). Tongue shape is also relevant to speech articulation, for example in retroflex consonants, where the tip of the tongue is curved backward.
Intimacy
The tongue plays a role in physical intimacy and sexuality. The tongue is part of the erogenous zone of the mouth and can be used in intimate contact, as in the French kiss and in oral sex.
Clinical significance
Disease
A congenital disorder of the tongue is that of ankyloglossia also known as tongue-tie. The tongue is tied to the floor of the mouth by a very short and thickened frenulum and this affects speech, eating, and swallowing.
The tongue is prone to several pathologies including glossitis and other inflammations such as geographic tongue, and median rhomboid glossitis; burning mouth syndrome, oral hairy leukoplakia, oral candidiasis (thrush), black hairy tongue, bifid tongue (due to failure in fusion of two lingual swellings of first pharyngeal arch) and fissured tongue.
There are several types of oral cancer that mainly affect the tongue. Mostly these are squamous cell carcinomas.
Food debris, desquamated epithelial cells and bacteria often form a visible tongue coating. This coating has been identified as a major factor contributing to bad breath (halitosis), which can be managed by using a tongue cleaner.
Medication delivery
The sublingual region underneath the front of the tongue is an ideal location for the administration of certain medications into the body. The oral mucosa is very thin underneath the tongue, and is underlain by a plexus of veins. The sublingual route takes advantage of the highly vascular quality of the oral cavity, and allows for the speedy application of medication into the cardiovascular system, bypassing the gastrointestinal tract. This is the only convenient and efficacious route of administration (apart from Intravenous therapy) of nitroglycerin to a patient suffering chest pain from angina pectoris.
Other animals
The muscles of the tongue evolved in amphibians from occipital somites. Most amphibians show a proper tongue after their metamorphosis. As a consequence, most tetrapod animals—amphibians, reptiles, birds, and mammals—have tongues (the frog family of pipids lack tongue). In mammals such as dogs and cats, the tongue is often used to clean the fur and body by licking. The tongues of these species have a very rough texture, which allows them to remove oils and parasites. Some dogs have a tendency to consistently lick a part of their foreleg, which can result in a skin condition known as a lick granuloma. A dog's tongue also acts as a heat regulator. As a dog increases its exercise the tongue will increase in size due to greater blood flow. The tongue hangs out of the dog's mouth and the moisture on the tongue will work to cool the bloodflow.
Some animals have tongues that are specially adapted for catching prey. For example, chameleons, frogs, pangolins and anteaters have prehensile tongues.
Other animals may have organs that are analogous to tongues, such as a butterfly's proboscis or a radula on a mollusc, but these are not homologous with the tongues found in vertebrates and often have little resemblance in function. For example, butterflies do not lick with their proboscides; they suck through them, and the proboscis is not a single organ, but two jaws held together to form a tube. Many species of fish have small folds at the base of their mouths that might informally be called tongues, but they lack a muscular structure like the true tongues found in most tetrapods.
Society and culture
Figures of speech
The tongue can serve as a metonym for language. For example, the New Testament of the Bible, in the Book of Acts of the Apostles, Jesus' disciples on the Day of Pentecost received a type of spiritual gift: "there appeared unto them cloven tongues like as of fire, and it sat upon each of them. And they were all filled with the Holy Ghost, and began to speak with other tongues ....", which amazed the crowd of Jewish people in Jerusalem, who were from various parts of the Roman Empire but could now understand what was being preached. The phrase mother tongue is used as a child's first language. Many languages have the same word for "tongue" and "language", as did the English language before the Middle Ages.
A common temporary failure in word retrieval from memory is referred to as the tip-of-the-tongue phenomenon. The expression tongue in cheek refers to a statement that is not to be taken entirely seriously – something said or done with subtle ironic or sarcastic humour. A tongue twister is a phrase very difficult to pronounce. Aside from being a medical condition, "tongue-tied" means being unable to say what you want due to confusion or restriction. The phrase "cat got your tongue" refers to when a person is speechless. To "bite one's tongue" is a phrase which describes holding back an opinion to avoid causing offence. A "slip of the tongue" refers to an unintentional utterance, such as a Freudian slip. The "gift of tongues" refers to when one is uncommonly gifted to be able to speak in a foreign language, often as a type of spiritual gift. Speaking in tongues is a common phrase used to describe glossolalia, which is to make smooth, language-resembling sounds that is no true spoken language itself. A deceptive person is said to have a forked tongue, and a smooth-talking person is said to have a .
Gestures
Sticking one's tongue out at someone is considered a childish gesture of rudeness or defiance in many countries; the act may also have sexual connotations, depending on the way in which it is done. However, in Tibet it is considered a greeting. In 2009, a farmer from Fabriano, Italy, was convicted and fined by Italy's highest court for sticking his tongue out at a neighbor with whom he had been arguing - proof of the affront had been captured with a cell-phone camera.
Body art
Tongue piercing and splitting have become more common in western countries in recent decades. One study found that one-fifth of young adults in Israel had at least one type of oral piercing, most commonly the tongue.
Representational art
Protruding tongues appear in the art of several Polynesian cultures.
As food
The tongues of some animals are consumed and sometimes prized as delicacies. Hot-tongue sandwiches frequently appear on menus in kosher delicatessens in America. Taco de lengua (lengua being Spanish for tongue) is a taco filled with beef tongue, and is especially popular in Mexican cuisine. As part of Colombian gastronomy, Tongue in Sauce (Lengua en Salsa) is a dish prepared by frying the tongue and adding tomato sauce, onions and salt. Tongue can also be prepared as birria. Pig and beef tongue are consumed in Chinese cuisine. Duck tongues are sometimes employed in Sichuan dishes, while lamb's tongue is occasionally employed in Continental and contemporary American cooking. Fried cod "tongue" is a relatively common part of fish meals in Norway and in Newfoundland. In Argentina and Uruguay cow tongue is cooked and served in vinegar (lengua a la vinagreta). In the Czech Republic and in Poland, a pork tongue is considered a delicacy, and there are many ways of preparing it. In Eastern Slavic countries, pork and beef tongues are commonly consumed, boiled and garnished with horseradish or jellied; beef tongues fetch a significantly higher price and are considered more of a delicacy. In Alaska, cow tongues are among the more common. Both cow and moose tongues are popular toppings on open-top-sandwiches in Norway, the latter usually amongst hunters.
Tongues of seals and whales have been eaten, sometimes in large quantities, by sealers and whalers, and in various times and places have been sold for food on shore.
Gallery
| Biology and health sciences | Digestive system | null |
56098 | https://en.wikipedia.org/wiki/Monte%20Carlo%20method | Monte Carlo method | Monte Carlo methods, or Monte Carlo experiments, are a broad class of computational algorithms that rely on repeated random sampling to obtain numerical results. The underlying concept is to use randomness to solve problems that might be deterministic in principle. The name comes from the Monte Carlo Casino in Monaco, where the primary developer of the method, mathematician Stanisław Ulam, was inspired by his uncle's gambling habits.
Monte Carlo methods are mainly used in three distinct problem classes: optimization, numerical integration, and generating draws from a probability distribution. They can also be used to model phenomena with significant uncertainty in inputs, such as calculating the risk of a nuclear power plant failure. Monte Carlo methods are often implemented using computer simulations, and they can provide approximate solutions to problems that are otherwise intractable or too complex to analyze mathematically.
Monte Carlo methods are widely used in various fields of science, engineering, and mathematics, such as physics, chemistry, biology, statistics, artificial intelligence, finance, and cryptography. They have also been applied to social sciences, such as sociology, psychology, and political science. Monte Carlo methods have been recognized as one of the most important and influential ideas of the 20th century, and they have enabled many scientific and technological breakthroughs.
Monte Carlo methods also have some limitations and challenges, such as the trade-off between accuracy and computational cost, the curse of dimensionality, the reliability of random number generators, and the verification and validation of the results.
Overview
Monte Carlo methods vary, but tend to follow a particular pattern:
Define a domain of possible inputs
Generate inputs randomly from a probability distribution over the domain
Perform a deterministic computation of the outputs
Aggregate the results
For example, consider a quadrant (circular sector) inscribed in a unit square. Given that the ratio of their areas is , the value of can be approximated using a Monte Carlo method:
Draw a square, then inscribe a quadrant within it
Uniformly scatter a given number of points over the square
Count the number of points inside the quadrant, i.e. having a distance from the origin of less than 1
The ratio of the inside-count and the total-sample-count is an estimate of the ratio of the two areas, . Multiply the result by 4 to estimate .
In this procedure the domain of inputs is the square that circumscribes the quadrant. One can generate random inputs by scattering grains over the square then perform a computation on each input (test whether it falls within the quadrant). Aggregating the results yields our final result, the approximation of .
There are two important considerations:
If the points are not uniformly distributed, then the approximation will be poor.
The approximation is generally poor if only a few points are randomly placed in the whole square. On average, the approximation improves as more points are placed.
Uses of Monte Carlo methods require large amounts of random numbers, and their use benefitted greatly from pseudorandom number generators, which are far quicker to use than the tables of random numbers that had been previously used for statistical sampling.
Application
Monte Carlo methods are often used in physical and mathematical problems and are most useful when it is difficult or impossible to use other approaches. Monte Carlo methods are mainly used in three problem classes: optimization, numerical integration, and generating draws from a probability distribution.
In physics-related problems, Monte Carlo methods are useful for simulating systems with many coupled degrees of freedom, such as fluids, disordered materials, strongly coupled solids, and cellular structures (see cellular Potts model, interacting particle systems, McKean–Vlasov processes, kinetic models of gases).
Other examples include modeling phenomena with significant uncertainty in inputs such as the calculation of risk in business and, in mathematics, evaluation of multidimensional definite integrals with complicated boundary conditions. In application to systems engineering problems (space, oil exploration, aircraft design, etc.), Monte Carlo–based predictions of failure, cost overruns and schedule overruns are routinely better than human intuition or alternative "soft" methods.
In principle, Monte Carlo methods can be used to solve any problem having a probabilistic interpretation. By the law of large numbers, integrals described by the expected value of some random variable can be approximated by taking the empirical mean ( the 'sample mean') of independent samples of the variable. When the probability distribution of the variable is parameterized, mathematicians often use a Markov chain Monte Carlo (MCMC) sampler. The central idea is to design a judicious Markov chain model with a prescribed stationary probability distribution. That is, in the limit, the samples being generated by the MCMC method will be samples from the desired (target) distribution. By the ergodic theorem, the stationary distribution is approximated by the empirical measures of the random states of the MCMC sampler.
In other problems, the objective is generating draws from a sequence of probability distributions satisfying a nonlinear evolution equation. These flows of probability distributions can always be interpreted as the distributions of the random states of a Markov process whose transition probabilities depend on the distributions of the current random states (see McKean–Vlasov processes, nonlinear filtering equation). In other instances, a flow of probability distributions with an increasing level of sampling complexity arise (path spaces models with an increasing time horizon, Boltzmann–Gibbs measures associated with decreasing temperature parameters, and many others). These models can also be seen as the evolution of the law of the random states of a nonlinear Markov chain. A natural way to simulate these sophisticated nonlinear Markov processes is to sample multiple copies of the process, replacing in the evolution equation the unknown distributions of the random states by the sampled empirical measures. In contrast with traditional Monte Carlo and MCMC methodologies, these mean-field particle techniques rely on sequential interacting samples. The terminology mean field reflects the fact that each of the samples ( particles, individuals, walkers, agents, creatures, or phenotypes) interacts with the empirical measures of the process. When the size of the system tends to infinity, these random empirical measures converge to the deterministic distribution of the random states of the nonlinear Markov chain, so that the statistical interaction between particles vanishes.
Simple Monte Carlo
Suppose one wants to know the expected value μ of a population (and knows that μ exists), but does not have a formula available to compute it. The simple Monte Carlo method gives an estimate for μ by running n simulations and averaging the simulations’ results. It has no restrictions on the probability distribution of the inputs to the simulations, requiring only that the inputs are randomly generated and are independent of each other and that μ exists. A sufficiently large n will produce a value for m that is arbitrarily close to μ; more formally, it will be the case that, for any ε > 0, |μ – m| ≤ ε.
Typically, the algorithm to obtain m is
s = 0;
for i = 1 to n do
run the simulation for the ith time, giving result ri;
s = s + ri;
repeat
m = s / n;
An example
Suppose we want to know how many times we should expect to throw three eight-sided dice for the total of the dice throws to be at least T. We know the expected value exists. The dice throws are randomly distributed and independent of each other. So simple Monte Carlo is applicable:
s = 0;
for i = 1 to n do
throw the three dice until T is met or first exceeded; ri = the number of throws;
s = s + ri;
repeat
m = s / n;
If n is large enough, m will be within ε of μ for any ε > 0.
Determining a sufficiently large n
General formula
Let ε = |μ – m| > 0. Choose the desired confidence level – the percent chance that, when the Monte Carlo algorithm completes, m is indeed within ε of μ. Let z be the z-score corresponding to that confidence level.
Let s2 be the estimated variance, sometimes called the “sample” variance; it is the variance of the results obtained from a relatively small number k of “sample” simulations. Choose a k; Driels and Shin observe that “even for sample sizes an order of magnitude lower than the number required, the calculation of that number is quite stable."
The following algorithm computes s2 in one pass while minimizing the possibility that accumulated numerical error produces erroneous results:
s1 = 0;
run the simulation for the first time, producing result r1;
m1 = r1; //mi is the mean of the first i simulations
for i = 2 to k do
run the simulation for the ith time, producing result ri;
δi = ri - mi−1;
mi = mi-1 + (1/i)δi;
si = si-1 + ((i - 1)/i)(δi)2;
repeat
s2 = sk/(k - 1);
Note that, when the algorithm completes, mk is the mean of the k results.
n is sufficiently large when
If n ≤ k, then mk = m; sufficient sample simulations were done to ensure that mk is within ε of μ. If n > k, then n simulations can be run “from scratch,” or, since k simulations have already been done, one can just run n – k more simulations and add their results into those from the sample simulations:
s = mk * k;
for i = k + 1 to n do
run the simulation for the ith time, giving result ri;
s = s + ri;
m = s / n;
A formula when simulations' results are bounded
An alternate formula can be used in the special case where all simulation results are bounded above and below.
Choose a value for ε that is twice the maximum allowed difference between μ and m. Let 0 < δ < 100 be the desired confidence level, expressed as a percentage. Let every simulation result r1, r2, …ri, … rn be such that a ≤ ri ≤ b for finite a and b. To have confidence of at least δ that |μ – m| < ε/2, use a value for n such that
For example, if δ = 99%, then n ≥ 2(b – a )2ln(2/0.01)/ε2 ≈ 10.6(b – a )2/ε2.
Computational costs
Despite its conceptual and algorithmic simplicity, the computational cost associated with a Monte Carlo simulation can be staggeringly high. In general the method requires many samples to get a good approximation, which may incur an arbitrarily large total runtime if the processing time of a single sample is high. Although this is a severe limitation in very complex problems, the embarrassingly parallel nature of the algorithm allows this large cost to be reduced (perhaps to a feasible level) through parallel computing strategies in local processors, clusters, cloud computing, GPU, FPGA, etc.
History
Before the Monte Carlo method was developed, simulations tested a previously understood deterministic problem, and statistical sampling was used to estimate uncertainties in the simulations. Monte Carlo simulations invert this approach, solving deterministic problems using probabilistic metaheuristics (see simulated annealing).
An early variant of the Monte Carlo method was devised to solve the Buffon's needle problem, in which can be estimated by dropping needles on a floor made of parallel equidistant strips. In the 1930s, Enrico Fermi first experimented with the Monte Carlo method while studying neutron diffusion, but he did not publish this work.
In the late 1940s, Stanisław Ulam invented the modern version of the Markov Chain Monte Carlo method while he was working on nuclear weapons projects at the Los Alamos National Laboratory. In 1946, nuclear weapons physicists at Los Alamos were investigating neutron diffusion in the core of a nuclear weapon. Despite having most of the necessary data, such as the average distance a neutron would travel in a substance before it collided with an atomic nucleus and how much energy the neutron was likely to give off following a collision, the Los Alamos physicists were unable to solve the problem using conventional, deterministic mathematical methods. Ulam proposed using random experiments. He recounts his inspiration as follows:
Being secret, the work of von Neumann and Ulam required a code name. A colleague of von Neumann and Ulam, Nicholas Metropolis, suggested using the name Monte Carlo, which refers to the Monte Carlo Casino in Monaco where Ulam's uncle would borrow money from relatives to gamble.
Monte Carlo methods were central to the simulations required for further postwar development of nuclear weapons, including the design of the H-bomb, though severely limited by the computational tools at the time. Von Neumann, Nicholas Metropolis and others programmed the ENIAC computer to perform the first fully automated Monte Carlo calculations, of a fission weapon core, in the spring of 1948. In the 1950s Monte Carlo methods were used at Los Alamos for the development of the hydrogen bomb, and became popularized in the fields of physics, physical chemistry, and operations research. The Rand Corporation and the U.S. Air Force were two of the major organizations responsible for funding and disseminating information on Monte Carlo methods during this time, and they began to find a wide application in many different fields.
The theory of more sophisticated mean-field type particle Monte Carlo methods had certainly started by the mid-1960s, with the work of Henry P. McKean Jr. on Markov interpretations of a class of nonlinear parabolic partial differential equations arising in fluid mechanics. An earlier pioneering article by Theodore E. Harris and Herman Kahn, published in 1951, used mean-field genetic-type Monte Carlo methods for estimating particle transmission energies. Mean-field genetic type Monte Carlo methodologies are also used as heuristic natural search algorithms (a.k.a. metaheuristic) in evolutionary computing. The origins of these mean-field computational techniques can be traced to 1950 and 1954 with the work of Alan Turing on genetic type mutation-selection learning machines and the articles by Nils Aall Barricelli at the Institute for Advanced Study in Princeton, New Jersey.
Quantum Monte Carlo, and more specifically diffusion Monte Carlo methods can also be interpreted as a mean-field particle Monte Carlo approximation of Feynman–Kac path integrals. The origins of Quantum Monte Carlo methods are often attributed to Enrico Fermi and Robert Richtmyer who developed in 1948 a mean-field particle interpretation of neutron-chain reactions, but the first heuristic-like and genetic type particle algorithm (a.k.a. Resampled or Reconfiguration Monte Carlo methods) for estimating ground state energies of quantum systems (in reduced matrix models) is due to Jack H. Hetherington in 1984. In molecular chemistry, the use of genetic heuristic-like particle methodologies (a.k.a. pruning and enrichment strategies) can be traced back to 1955 with the seminal work of Marshall N. Rosenbluth and Arianna W. Rosenbluth.
The use of Sequential Monte Carlo in advanced signal processing and Bayesian inference is more recent. It was in 1993, that Gordon et al., published in their seminal work the first application of a Monte Carlo resampling algorithm in Bayesian statistical inference. The authors named their algorithm 'the bootstrap filter', and demonstrated that compared to other filtering methods, their bootstrap algorithm does not require any assumption about that state-space or the noise of the system. Another pioneering article in this field was Genshiro Kitagawa's, on a related "Monte Carlo filter", and the ones by Pierre Del Moral and Himilcon Carvalho, Pierre Del Moral, André Monin and Gérard Salut on particle filters published in the mid-1990s. Particle filters were also developed in signal processing in 1989–1992 by P. Del Moral, J. C. Noyer, G. Rigal, and G. Salut in the LAAS-CNRS in a series of restricted and classified research reports with STCAN (Service Technique des Constructions et Armes Navales), the IT company DIGILOG, and the LAAS-CNRS (the Laboratory for Analysis and Architecture of Systems) on radar/sonar and GPS signal processing problems. These Sequential Monte Carlo methodologies can be interpreted as an acceptance-rejection sampler equipped with an interacting recycling mechanism.
From 1950 to 1996, all the publications on Sequential Monte Carlo methodologies, including the pruning and resample Monte Carlo methods introduced in computational physics and molecular chemistry, present natural and heuristic-like algorithms applied to different situations without a single proof of their consistency, nor a discussion on the bias of the estimates and on genealogical and ancestral tree based algorithms. The mathematical foundations and the first rigorous analysis of these particle algorithms were written by Pierre Del Moral in 1996.
Branching type particle methodologies with varying population sizes were also developed in the end of the 1990s by Dan Crisan, Jessica Gaines and Terry Lyons, and by Dan Crisan, Pierre Del Moral and Terry Lyons. Further developments in this field were described in 1999 to 2001 by P. Del Moral, A. Guionnet and L. Miclo.
Definitions
There is no consensus on how Monte Carlo should be defined. For example, Ripley defines most probabilistic modeling as stochastic simulation, with Monte Carlo being reserved for Monte Carlo integration and Monte Carlo statistical tests. Sawilowsky distinguishes between a simulation, a Monte Carlo method, and a Monte Carlo simulation: a simulation is a fictitious representation of reality, a Monte Carlo method is a technique that can be used to solve a mathematical or statistical problem, and a Monte Carlo simulation uses repeated sampling to obtain the statistical properties of some phenomenon (or behavior).
Here are some examples:
Simulation: Drawing one pseudo-random uniform variable from the interval [0,1] can be used to simulate the tossing of a coin: If the value is less than or equal to 0.50 designate the outcome as heads, but if the value is greater than 0.50 designate the outcome as tails. This is a simulation, but not a Monte Carlo simulation.
Monte Carlo method: Pouring out a box of coins on a table, and then computing the ratio of coins that land heads versus tails is a Monte Carlo method of determining the behavior of repeated coin tosses, but it is not a simulation.
Monte Carlo simulation: Drawing a large number of pseudo-random uniform variables from the interval [0,1] at one time, or once at many different times, and assigning values less than or equal to 0.50 as heads and greater than 0.50 as tails, is a Monte Carlo simulation of the behavior of repeatedly tossing a coin.
Kalos and Whitlock point out that such distinctions are not always easy to maintain. For example, the emission of radiation from atoms is a natural stochastic process. It can be simulated directly, or its average behavior can be described by stochastic equations that can themselves be solved using Monte Carlo methods. "Indeed, the same computer code can be viewed simultaneously as a 'natural simulation' or as a solution of the equations by natural sampling."
Convergence of the Monte Carlo simulation can be checked with the Gelman-Rubin statistic.
Monte Carlo and random numbers
The main idea behind this method is that the results are computed based on repeated random sampling and statistical analysis. The Monte Carlo simulation is, in fact, random experimentations, in the case that, the results of these experiments are not well known.
Monte Carlo simulations are typically characterized by many unknown parameters, many of which are difficult to obtain experimentally. Monte Carlo simulation methods do not always require truly random numbers to be useful (although, for some applications such as primality testing, unpredictability is vital). Many of the most useful techniques use deterministic, pseudorandom sequences, making it easy to test and re-run simulations. The only quality usually necessary to make good simulations is for the pseudo-random sequence to appear "random enough" in a certain sense.
What this means depends on the application, but typically they should pass a series of statistical tests. Testing that the numbers are uniformly distributed or follow another desired distribution when a large enough number of elements of the sequence are considered is one of the simplest and most common ones. Weak correlations between successive samples are also often desirable/necessary.
Sawilowsky lists the characteristics of a high-quality Monte Carlo simulation:
the (pseudo-random) number generator has certain characteristics (e.g. a long "period" before the sequence repeats)
the (pseudo-random) number generator produces values that pass tests for randomness
there are enough samples to ensure accurate results
the proper sampling technique is used
the algorithm used is valid for what is being modeled
it simulates the phenomenon in question.
Pseudo-random number sampling algorithms are used to transform uniformly distributed pseudo-random numbers into numbers that are distributed according to a given probability distribution.
Low-discrepancy sequences are often used instead of random sampling from a space as they ensure even coverage and normally have a faster order of convergence than Monte Carlo simulations using random or pseudorandom sequences. Methods based on their use are called quasi-Monte Carlo methods.
In an effort to assess the impact of random number quality on Monte Carlo simulation outcomes, astrophysical researchers tested cryptographically secure pseudorandom numbers generated via Intel's RDRAND instruction set, as compared to those derived from algorithms, like the Mersenne Twister, in Monte Carlo simulations of radio flares from brown dwarfs. No statistically significant difference was found between models generated with typical pseudorandom number generators and RDRAND for trials consisting of the generation of 107 random numbers.
Monte Carlo simulation versus "what if" scenarios
There are ways of using probabilities that are definitely not Monte Carlo simulations – for example, deterministic modeling using single-point estimates. Each uncertain variable within a model is assigned a "best guess" estimate. Scenarios (such as best, worst, or most likely case) for each input variable are chosen and the results recorded.
By contrast, Monte Carlo simulations sample from a probability distribution for each variable to produce hundreds or thousands of possible outcomes. The results are analyzed to get probabilities of different outcomes occurring. For example, a comparison of a spreadsheet cost construction model run using traditional "what if" scenarios, and then running the comparison again with Monte Carlo simulation and triangular probability distributions shows that the Monte Carlo analysis has a narrower range than the "what if" analysis. This is because the "what if" analysis gives equal weight to all scenarios (see quantifying uncertainty in corporate finance), while the Monte Carlo method hardly samples in the very low probability regions. The samples in such regions are called "rare events".
Applications
Monte Carlo methods are especially useful for simulating phenomena with significant uncertainty in inputs and systems with many coupled degrees of freedom. Areas of application include:
Physical sciences
Monte Carlo methods are very important in computational physics, physical chemistry, and related applied fields, and have diverse applications from complicated quantum chromodynamics calculations to designing heat shields and aerodynamic forms as well as in modeling radiation transport for radiation dosimetry calculations. In statistical physics, Monte Carlo molecular modeling is an alternative to computational molecular dynamics, and Monte Carlo methods are used to compute statistical field theories of simple particle and polymer systems. Quantum Monte Carlo methods solve the many-body problem for quantum systems. In radiation materials science, the binary collision approximation for simulating ion implantation is usually based on a Monte Carlo approach to select the next colliding atom. In experimental particle physics, Monte Carlo methods are used for designing detectors, understanding their behavior and comparing experimental data to theory. In astrophysics, they are used in such diverse manners as to model both galaxy evolution and microwave radiation transmission through a rough planetary surface. Monte Carlo methods are also used in the ensemble models that form the basis of modern weather forecasting.
Engineering
Monte Carlo methods are widely used in engineering for sensitivity analysis and quantitative probabilistic analysis in process design. The need arises from the interactive, co-linear and non-linear behavior of typical process simulations. For example,
In microelectronics engineering, Monte Carlo methods are applied to analyze correlated and uncorrelated variations in analog and digital integrated circuits.
In geostatistics and geometallurgy, Monte Carlo methods underpin the design of mineral processing flowsheets and contribute to quantitative risk analysis.
In fluid dynamics, in particular rarefied gas dynamics, where the Boltzmann equation is solved for finite Knudsen number fluid flows using the direct simulation Monte Carlo method in combination with highly efficient computational algorithms.
In autonomous robotics, Monte Carlo localization can determine the position of a robot. It is often applied to stochastic filters such as the Kalman filter or particle filter that forms the heart of the SLAM (simultaneous localization and mapping) algorithm.
In telecommunications, when planning a wireless network, the design must be proven to work for a wide variety of scenarios that depend mainly on the number of users, their locations and the services they want to use. Monte Carlo methods are typically used to generate these users and their states. The network performance is then evaluated and, if results are not satisfactory, the network design goes through an optimization process.
In reliability engineering, Monte Carlo simulation is used to compute system-level response given the component-level response.
In signal processing and Bayesian inference, particle filters and sequential Monte Carlo techniques are a class of mean-field particle methods for sampling and computing the posterior distribution of a signal process given some noisy and partial observations using interacting empirical measures.
Climate change and radiative forcing
The Intergovernmental Panel on Climate Change relies on Monte Carlo methods in probability density function analysis of radiative forcing.
Computational biology
Monte Carlo methods are used in various fields of computational biology, for example for Bayesian inference in phylogeny, or for studying biological systems such as genomes, proteins, or membranes.
The systems can be studied in the coarse-grained or ab initio frameworks depending on the desired accuracy.
Computer simulations allow monitoring of the local environment of a particular molecule to see if some chemical reaction is happening for instance. In cases where it is not feasible to conduct a physical experiment, thought experiments can be conducted (for instance: breaking bonds, introducing impurities at specific sites, changing the local/global structure, or introducing external fields).
Computer graphics
Path tracing, occasionally referred to as Monte Carlo ray tracing, renders a 3D scene by randomly tracing samples of possible light paths. Repeated sampling of any given pixel will eventually cause the average of the samples to converge on the correct solution of the rendering equation, making it one of the most physically accurate 3D graphics rendering methods in existence.
Applied statistics
The standards for Monte Carlo experiments in statistics were set by Sawilowsky. In applied statistics, Monte Carlo methods may be used for at least four purposes:
To compare competing statistics for small samples under realistic data conditions. Although type I error and power properties of statistics can be calculated for data drawn from classical theoretical distributions (e.g., normal curve, Cauchy distribution) for asymptotic conditions (i. e, infinite sample size and infinitesimally small treatment effect), real data often do not have such distributions.
To provide implementations of hypothesis tests that are more efficient than exact tests such as permutation tests (which are often impossible to compute) while being more accurate than critical values for asymptotic distributions.
To provide a random sample from the posterior distribution in Bayesian inference. This sample then approximates and summarizes all the essential features of the posterior.
To provide efficient random estimates of the Hessian matrix of the negative log-likelihood function that may be averaged to form an estimate of the Fisher information matrix.
Monte Carlo methods are also a compromise between approximate randomization and permutation tests. An approximate randomization test is based on a specified subset of all permutations (which entails potentially enormous housekeeping of which permutations have been considered). The Monte Carlo approach is based on a specified number of randomly drawn permutations (exchanging a minor loss in precision if a permutation is drawn twice—or more frequently—for the efficiency of not having to track which permutations have already been selected).
Artificial intelligence for games
Monte Carlo methods have been developed into a technique called Monte-Carlo tree search that is useful for searching for the best move in a game. Possible moves are organized in a search tree and many random simulations are used to estimate the long-term potential of each move. A black box simulator represents the opponent's moves.
The Monte Carlo tree search (MCTS) method has four steps:
Starting at root node of the tree, select optimal child nodes until a leaf node is reached.
Expand the leaf node and choose one of its children.
Play a simulated game starting with that node.
Use the results of that simulated game to update the node and its ancestors.
The net effect, over the course of many simulated games, is that the value of a node representing a move will go up or down, hopefully corresponding to whether or not that node represents a good move.
Monte Carlo Tree Search has been used successfully to play games such as Go, Tantrix, Battleship, Havannah, and Arimaa.
Design and visuals
Monte Carlo methods are also efficient in solving coupled integral differential equations of radiation fields and energy transport, and thus these methods have been used in global illumination computations that produce photo-realistic images of virtual 3D models, with applications in video games, architecture, design, computer generated films, and cinematic special effects.
Search and rescue
The US Coast Guard utilizes Monte Carlo methods within its computer modeling software SAROPS in order to calculate the probable locations of vessels during search and rescue operations. Each simulation can generate as many as ten thousand data points that are randomly distributed based upon provided variables. Search patterns are then generated based upon extrapolations of these data in order to optimize the probability of containment (POC) and the probability of detection (POD), which together will equal an overall probability of success (POS). Ultimately this serves as a practical application of probability distribution in order to provide the swiftest and most expedient method of rescue, saving both lives and resources.
Finance and business
Monte Carlo simulation is commonly used to evaluate the risk and uncertainty that would affect the outcome of different decision options. Monte Carlo simulation allows the business risk analyst to incorporate the total effects of uncertainty in variables like sales volume, commodity and labor prices, interest and exchange rates, as well as the effect of distinct risk events like the cancellation of a contract or the change of a tax law.
Monte Carlo methods in finance are often used to evaluate investments in projects at a business unit or corporate level, or other financial valuations. They can be used to model project schedules, where simulations aggregate estimates for worst-case, best-case, and most likely durations for each task to determine outcomes for the overall project. Monte Carlo methods are also used in option pricing, default risk analysis. Additionally, they can be used to estimate the financial impact of medical interventions.
Law
A Monte Carlo approach was used for evaluating the potential value of a proposed program to help female petitioners in Wisconsin be successful in their applications for harassment and domestic abuse restraining orders. It was proposed to help women succeed in their petitions by providing them with greater advocacy thereby potentially reducing the risk of rape and physical assault. However, there were many variables in play that could not be estimated perfectly, including the effectiveness of restraining orders, the success rate of petitioners both with and without advocacy, and many others. The study ran trials that varied these variables to come up with an overall estimate of the success level of the proposed program as a whole.
Library science
Monte Carlo approach had also been used to simulate the number of book publications based on book genre in Malaysia. The Monte Carlo simulation utilized previous published National Book publication data and book's price according to book genre in the local market. The Monte Carlo results were used to determine what kind of book genre that Malaysians are fond of and was used to compare book publications between Malaysia and Japan.
Other
Nassim Nicholas Taleb writes about Monte Carlo generators in his 2001 book Fooled by Randomness as a real instance of the reverse Turing test: a human can be declared unintelligent if their writing cannot be told apart from a generated one.
Use in mathematics
In general, the Monte Carlo methods are used in mathematics to solve various problems by generating suitable random numbers (see also Random number generation) and observing that fraction of the numbers that obeys some property or properties. The method is useful for obtaining numerical solutions to problems too complicated to solve analytically. The most common application of the Monte Carlo method is Monte Carlo integration.
Integration
Deterministic numerical integration algorithms work well in a small number of dimensions, but encounter two problems when the functions have many variables. First, the number of function evaluations needed increases rapidly with the number of dimensions. For example, if 10 evaluations provide adequate accuracy in one dimension, then 10100 points are needed for 100 dimensions—far too many to be computed. This is called the curse of dimensionality. Second, the boundary of a multidimensional region may be very complicated, so it may not be feasible to reduce the problem to an iterated integral. 100 dimensions is by no means unusual, since in many physical problems, a "dimension" is equivalent to a degree of freedom.
Monte Carlo methods provide a way out of this exponential increase in computation time. As long as the function in question is reasonably well-behaved, it can be estimated by randomly selecting points in 100-dimensional space, and taking some kind of average of the function values at these points. By the central limit theorem, this method displays convergence—i.e., quadrupling the number of sampled points halves the error, regardless of the number of dimensions.
A refinement of this method, known as importance sampling in statistics, involves sampling the points randomly, but more frequently where the integrand is large. To do this precisely one would have to already know the integral, but one can approximate the integral by an integral of a similar function or use adaptive routines such as stratified sampling, recursive stratified sampling, adaptive umbrella sampling or the VEGAS algorithm.
A similar approach, the quasi-Monte Carlo method, uses low-discrepancy sequences. These sequences "fill" the area better and sample the most important points more frequently, so quasi-Monte Carlo methods can often converge on the integral more quickly.
Another class of methods for sampling points in a volume is to simulate random walks over it (Markov chain Monte Carlo). Such methods include the Metropolis–Hastings algorithm, Gibbs sampling, Wang and Landau algorithm, and interacting type MCMC methodologies such as the sequential Monte Carlo samplers.
Simulation and optimization
Another powerful and very popular application for random numbers in numerical simulation is in numerical optimization. The problem is to minimize (or maximize) functions of some vector that often has many dimensions. Many problems can be phrased in this way: for example, a computer chess program could be seen as trying to find the set of, say, 10 moves that produces the best evaluation function at the end. In the traveling salesman problem the goal is to minimize distance traveled. There are also applications to engineering design, such as multidisciplinary design optimization. It has been applied with quasi-one-dimensional models to solve particle dynamics problems by efficiently exploring large configuration space. Reference is a comprehensive review of many issues related to simulation and optimization.
The traveling salesman problem is what is called a conventional optimization problem. That is, all the facts (distances between each destination point) needed to determine the optimal path to follow are known with certainty and the goal is to run through the possible travel choices to come up with the one with the lowest total distance. If instead of the goal being to minimize the total distance traveled to visit each desired destination but rather to minimize the total time needed to reach each destination, this goes beyond conventional optimization since travel time is inherently uncertain (traffic jams, time of day, etc.). As a result, to determine the optimal path a different simulation is required: optimization to first understand the range of potential times it could take to go from one point to another (represented by a probability distribution in this case rather than a specific distance) and then optimize the travel decisions to identify the best path to follow taking that uncertainty into account.
Inverse problems
Probabilistic formulation of inverse problems leads to the definition of a probability distribution in the model space. This probability distribution combines prior information with new information obtained by measuring some observable parameters (data).
As, in the general case, the theory linking data with model parameters is nonlinear, the posterior probability in the model space may not be easy to describe (it may be multimodal, some moments may not be defined, etc.).
When analyzing an inverse problem, obtaining a maximum likelihood model is usually not sufficient, as normally information on the resolution power of the data is desired. In the general case many parameters are modeled, and an inspection of the marginal probability densities of interest may be impractical, or even useless. But it is possible to pseudorandomly generate a large collection of models according to the posterior probability distribution and to analyze and display the models in such a way that information on the relative likelihoods of model properties is conveyed to the spectator. This can be accomplished by means of an efficient Monte Carlo method, even in cases where no explicit formula for the a priori distribution is available.
The best-known importance sampling method, the Metropolis algorithm, can be generalized, and this gives a method that allows analysis of (possibly highly nonlinear) inverse problems with complex a priori information and data with an arbitrary noise distribution.
Philosophy
Popular exposition of the Monte Carlo Method was conducted by McCracken. The method's general philosophy was discussed by Elishakoff and Grüne-Yanoff and Weirich.
| Mathematics | Statistics | null |
56099 | https://en.wikipedia.org/wiki/Red%20dwarf | Red dwarf | A red dwarf is the smallest kind of star on the main sequence. Red dwarfs are by far the most common type of fusing star in the Milky Way, at least in the neighborhood of the Sun. However, due to their low luminosity, individual red dwarfs cannot be easily observed. From Earth, not one star that fits the stricter definitions of a red dwarf is visible to the naked eye. Proxima Centauri, the star nearest to the Sun, is a red dwarf, as are fifty of the sixty nearest stars. According to some estimates, red dwarfs make up three-quarters of the fusing stars in the Milky Way.
The coolest red dwarfs near the Sun have a surface temperature of about and the smallest have radii about 9% that of the Sun, with masses about 7.5% that of the Sun. These red dwarfs have spectral types of L0 to L2. There is some overlap with the properties of brown dwarfs, since the most massive brown dwarfs at lower metallicity can be as hot as and have late M spectral types.
Definitions and usage of the term "red dwarf" vary on how inclusive they are on the hotter and more massive end. One definition is synonymous with stellar M dwarfs (M-type main sequence stars), yielding a maximum temperature of and . One includes all stellar M-type main-sequence and all K-type main-sequence stars (K dwarf), yielding a maximum temperature of and . Some definitions include any stellar M dwarf and part of the K dwarf classification. Other definitions are also in use. Many of the coolest, lowest mass M dwarfs are expected to be brown dwarfs, not true stars, and so those would be excluded from any definition of red dwarf.
Stellar models indicate that red dwarfs less than are fully convective. Hence, the helium produced by the thermonuclear fusion of hydrogen is constantly remixed throughout the star, avoiding helium buildup at the core, thereby prolonging the period of fusion. Low-mass red dwarfs therefore develop very slowly, maintaining a constant luminosity and spectral type for trillions of years, until their fuel is depleted. Because of the comparatively short age of the universe, no red dwarfs yet exist at advanced stages of evolution.
Definition
The term "red dwarf" when used to refer to a star does not have a strict definition. One of the earliest uses of the term was in 1915, used simply to contrast "red" dwarf stars from hotter "blue" dwarf stars. It became established use, although the definition remained vague. In terms of which spectral types qualify as red dwarfs, different researchers picked different limits, for example K8–M5 or "later than K5". Dwarf M star, abbreviated dM, was also used, but sometimes it also included stars of spectral type K.
In modern usage, the definition of a red dwarf still varies. When explicitly defined, it typically includes late K- and early to mid-M-class stars, but in many cases it is restricted just to M-class stars. In some cases all K stars are included as red dwarfs, and occasionally even earlier stars.
The most recent surveys place the coolest true main-sequence stars into spectral types L2 or L3. At the same time, many objects cooler than about M6 or M7 are brown dwarfs, insufficiently massive to sustain hydrogen-1 fusion. This gives a significant overlap in spectral types for red and brown dwarfs. Objects in that spectral range can be difficult to categorize.
Description and characteristics
Red dwarfs are very-low-mass stars. As a result, they have relatively low pressures, a low fusion rate, and hence, a low temperature. The energy generated is the product of nuclear fusion of hydrogen into helium by way of the proton–proton (PP) chain mechanism. Hence, these stars emit relatively little light, sometimes as little as that of the Sun, although this would still imply a power output on the order of 1022 watts (10 trillion gigawatts or 10 ZW). Even the largest red dwarfs (for example HD 179930, HIP 12961 and Lacaille 8760) have only about 10% of the Sun's luminosity. In general, red dwarfs less than transport energy from the core to the surface by convection. Convection occurs because of opacity of the interior, which has a high density compared to the temperature. As a result, energy transfer by radiation is decreased, and instead convection is the main form of energy transport to the surface of the star. Above this mass, a red dwarf will have a region around its core where convection does not occur.
Because low-mass red dwarfs are fully convective, helium does not accumulate at the core, and compared to larger stars such as the Sun, they can burn a larger proportion of their hydrogen before leaving the main sequence. As a result, red dwarfs have estimated lifespans far longer than the present age of the universe, and stars less than have not had time to leave the main sequence. The lower the mass of a red dwarf, the longer the lifespan. It is believed that the lifespan of these stars exceeds the expected 10-billion-year lifespan of the Sun by the third or fourth power of the ratio of the solar mass to their masses; thus, a red dwarf may continue burning for 10 trillion years. As the proportion of hydrogen in a red dwarf is consumed, the rate of fusion declines and the core starts to contract. The gravitational energy released by this size reduction is converted into heat, which is carried throughout the star by convection.
According to computer simulations, the minimum mass a red dwarf must have to eventually evolve into a red giant is ; less massive objects, as they age, would increase their surface temperatures and luminosities becoming blue dwarfs and finally white dwarfs.
The less massive the star, the longer this evolutionary process takes. A red dwarf (approximately the mass of the nearby Barnard's Star) would stay on the main sequence for 2.5 trillion years, followed by five billion years as a blue dwarf, during which the star would have one third of the Sun's luminosity () and a surface temperature of 6,500–8,500 kelvins.
The fact that red dwarfs and other low-mass stars still remain on the main sequence when more massive stars have moved off the main sequence allows the age of star clusters to be estimated by finding the mass at which the stars move off the main sequence. This provides a lower limit to the age of the Universe and also allows formation timescales to be placed upon the structures within the Milky Way, such as the Galactic halo and Galactic disk.
All observed red dwarfs contain "metals", which in astronomy are elements heavier than hydrogen and helium. The Big Bang model predicts that the first generation of stars should have only hydrogen, helium, and trace amounts of lithium, and hence would be of low metallicity. With their extreme lifespans, any red dwarfs that were a part of that first generation (population III stars) should still exist today. Low-metallicity red dwarfs, however, are rare. The accepted model for the chemical evolution of the universe anticipates such a scarcity of metal-poor dwarf stars because only giant stars are thought to have formed in the metal-poor environment of the early universe. As giant stars end their short lives in supernova explosions, they spew out the heavier elements needed to form smaller stars. Therefore, dwarfs became more common as the universe aged and became enriched in metals. While the basic scarcity of ancient metal-poor red dwarfs is expected, observations have detected even fewer than predicted. The sheer difficulty of detecting objects as dim as red dwarfs was thought to account for this discrepancy, but improved detection methods have only confirmed the discrepancy.
The boundary between the least massive red dwarfs and the most massive brown dwarfs depends strongly on the metallicity. At solar metallicity the boundary occurs at about , while at zero metallicity the boundary is around . At solar metallicity, the least massive red dwarfs theoretically have temperatures around , while measurements of red dwarfs in the solar neighbourhood suggest the coolest stars have temperatures of about and spectral classes of about L2. Theory predicts that the coolest red dwarfs at zero metallicity would have temperatures of about . The least massive red dwarfs have radii of about , while both more massive red dwarfs and less massive brown dwarfs are larger.
Spectral standard stars
The spectral standards for M type stars have changed slightly over the years, but settled down somewhat since the early 1990s. Part of this is due to the fact that even the nearest red dwarfs are fairly faint, and their colors do not register well on photographic emulsions used in the early to mid 20th century. The study of mid- to late-M dwarfs has significantly advanced only in the past few decades, primarily due to development of new astrographic and spectroscopic techniques, dispensing with photographic plates and progressing to charged-couple devices (CCDs) and infrared-sensitive arrays.
The revised Yerkes Atlas system (Johnson & Morgan, 1953) listed only two M type spectral standard stars: HD 147379 (M0V)
and HD 95735/Lalande 21185 (M2V). While HD 147379 was not considered a standard by expert classifiers in later compendia of standards, Lalande 21185 is still a primary standard for M2V. Robert Garrison does not list any "anchor" standards among the red dwarfs, but Lalande 21185 has survived as a M2V standard through many compendia. The review on MK classification by Morgan & Keenan (1973) did not contain red dwarf standards.
In the mid-1970s, red dwarf standard stars were published by Keenan & McNeil (1976) and Boeshaar (1976), but there was little agreement among the standards. As later cooler stars were identified through the 1980s, it was clear that an overhaul of the red dwarf standards was needed. Building primarily upon the Boeshaar standards, a group at Steward Observatory (Kirkpatrick, Henry, & McCarthy, 1991) filled in the spectral sequence from K5V to M9V. It is these M type dwarf standard stars which have largely survived as the main standards to the modern day. There have been negligible changes in the red dwarf spectral sequence since 1991. Additional red dwarf standards were compiled by Henry et al. (2002), and D. Kirkpatrick has recently
reviewed the classification of red dwarfs and standard stars in Gray & Corbally's 2009 monograph. The M dwarf primary spectral standards are: GJ 270 (M0V), GJ 229A (M1V), Lalande 21185 (M2V), Gliese 581 (M3V), Gliese 402 (M4V), GJ 51 (M5V), Wolf 359 (M6V), van Biesbroeck 8 (M7V), VB 10 (M8V), LHS 2924 (M9V).
Planets
Many red dwarfs are orbited by exoplanets, but large Jupiter-sized planets are comparatively rare. Doppler surveys of a wide variety of stars indicate about 1 in 6 stars with twice the mass of the Sun are orbited by one or more of Jupiter-sized planets, versus 1 in 16 for Sun-like stars and the frequency of close-in giant planets (Jupiter size or larger) orbiting red dwarfs is only 1 in 40. On the other hand, microlensing surveys indicate that long-orbital-period Neptune-mass planets are found around one in three red dwarfs. Observations with HARPS further indicate 40% of red dwarfs have a "super-Earth" class planet orbiting in the habitable zone where liquid water can exist on the surface. Computer simulations of the formation of planets around low-mass stars predict that Earth-sized planets are most abundant, but more than 90% of the simulated planets are at least 10% water by mass, suggesting that many Earth-sized planets orbiting red dwarf stars are covered in deep oceans.
At least four and possibly up to six exoplanets were discovered orbiting within the Gliese 581 planetary system between 2005 and 2010. One planet has about the mass of Neptune, or 16 Earth masses (). It orbits just from its star, and is estimated to have a surface temperature of , despite the dimness of its star. In 2006, an even smaller exoplanet (only ) was found orbiting the red dwarf OGLE-2005-BLG-390L; it lies from the star and its surface temperature is .
In 2007, a new, potentially habitable exoplanet, , was found, orbiting Gliese 581. The minimum mass estimated by its discoverers (a team led by Stephane Udry) is . The discoverers estimate its radius to be 1.5 times that of Earth (). Since then Gliese 581d, which is also potentially habitable, was discovered.
Gliese 581c and d are within the habitable zone of the host star, and are two of the most likely candidates for habitability of any exoplanets discovered so far. Gliese 581g, detected September 2010, has a near-circular orbit in the middle of the star's habitable zone. However, the planet's existence is contested.
On 23 February 2017 NASA announced the discovery of seven Earth-sized planets orbiting the red dwarf star TRAPPIST-1 approximately 39 light-years away in the constellation Aquarius. The planets were discovered through the transit method, meaning we have mass and radius information for all of them. TRAPPIST-1e, f, and g appear to be within the habitable zone and may have liquid water on the surface.
Habitability
Modern evidence suggests that planets in red dwarf systems are extremely unlikely to be habitable. In spite of their great numbers and long lifespans, there are several factors which may make life difficult on planets around a red dwarf. First, planets in the habitable zone of a red dwarf would be so close to the parent star that they would likely be tidally locked. For a nearly circular orbit, this would mean that one side would be in perpetual daylight and the other in eternal night. This could create enormous temperature variations from one side of the planet to the other. Such conditions would appear to make it difficult for forms of life similar to those on Earth to evolve. And it appears there is a great problem with the atmosphere of such tidally locked planets: the perpetual night zone would be cold enough to freeze the main gases of their atmospheres, leaving the daylight zone bare and dry. On the other hand, though, a theory proposes that either a thick atmosphere or planetary ocean could potentially circulate heat around such a planet.
Variability in stellar energy output may also have negative impacts on the development of life. Red dwarfs are often flare stars, which can emit gigantic flares, doubling their brightness in minutes. This variability makes it difficult for life to develop and persist near a red dwarf. While it may be possible for a planet orbiting close to a red dwarf to keep its atmosphere even if the star flares, more-recent research suggests that these stars may be the source of constant high-energy flares and very large magnetic fields, diminishing the possibility of life as we know it.
| Physical sciences | Stellar astronomy | Astronomy |
56106 | https://en.wikipedia.org/wiki/Wildfire | Wildfire | A wildfire, forest fire, or a bushfire is an unplanned, uncontrolled and unpredictable fire in an area of combustible vegetation. Depending on the type of vegetation present, a wildfire may be more specifically identified as a bushfire (in Australia), desert fire, grass fire, hill fire, peat fire, prairie fire, vegetation fire, or veld fire. Some natural forest ecosystems depend on wildfire. Modern forest management often engages in prescribed burns to mitigate fire risk and promote natural forest cycles. However, controlled burns can turn into wildfires by mistake.
Wildfires can be classified by cause of ignition, physical properties, combustible material present, and the effect of weather on the fire. Wildfire severity results from a combination of factors such as available fuels, physical setting, and weather. Climatic cycles with wet periods that create substantial fuels, followed by drought and heat, often precede severe wildfires. These cycles have been intensified by climate change.
Wildfires are a common type of disaster in some regions, including Siberia (Russia), California, Washington, Oregon, Texas, Florida, (United States), British Columbia (Canada), and Australia. Areas with Mediterranean climates or in the taiga biome are particularly susceptible. Wildfires can severely impact humans and their settlements. Effects include for example the direct health impacts of smoke and fire, as well as destruction of property (especially in wildland–urban interfaces), and economic losses. There is also the potential for contamination of water and soil.
At a global level, human practices have made the impacts of wildfire worse, with a doubling in land area burned by wildfires compared to natural levels. Humans have impacted wildfire through climate change (e.g. more intense heat waves and droughts), land-use change, and wildfire suppression. The carbon released from wildfires can add to carbon dioxide concentrations in the atmosphere and thus contribute to the greenhouse effect. This creates a climate change feedback.
Naturally occurring wildfires can have beneficial effects on those ecosystems that have evolved with fire. In fact, many plant species depend on the effects of fire for growth and reproduction.
Ignition
The ignition of a fire takes place through either natural causes or human activity (deliberate or not).
Natural causes
Natural occurrences that can ignite wildfires without the involvement of humans include lightning, volcanic eruptions, sparks from rock falls, and spontaneous combustions.
Human activity
Sources of human-caused fire may include arson, accidental ignition, or the uncontrolled use of fire in land-clearing and agriculture such as the slash-and-burn farming in Southeast Asia. In the tropics, farmers often practice the slash-and-burn method of clearing fields during the dry season.
In middle latitudes, the most common human causes of wildfires are equipment generating sparks (chainsaws, grinders, mowers, etc.), overhead power lines, and arson.
Arson may account for over 20% of human caused fires. However, in the 2019–20 Australian bushfire season "an independent study found online bots and trolls exaggerating the role of arson in the fires." In the 2023 Canadian wildfires false claims of arson gained traction on social media; however, arson is generally not a main cause of wildfires in Canada. In California, generally 6–10% of wildfires annually are arson.
Coal seam fires burn in the thousands around the world, such as those in Burning Mountain, New South Wales; Centralia, Pennsylvania; and several coal-sustained fires in China. They can also flare up unexpectedly and ignite nearby flammable material.
Spread
The spread of wildfires varies based on the flammable material present, its vertical arrangement and moisture content, and weather conditions. Fuel arrangement and density is governed in part by topography, as land shape determines factors such as available sunlight and water for plant growth. Overall, fire types can be generally characterized by their fuels as follows:
Ground fires are fed by subterranean roots, duff on the forest floor, and other buried organic matter. Ground fires typically burn by smoldering, and can burn slowly for days to months, such as peat fires in Kalimantan and Eastern Sumatra, Indonesia, which resulted from a riceland creation project that unintentionally drained and dried the peat.
Crawling or surface fires are fueled by low-lying vegetative matter on the forest floor such as leaf and timber litter, debris, grass, and low-lying shrubbery. This kind of fire often burns at a relatively lower temperature than crown fires (less than ) and may spread at slow rate, though steep slopes and wind can accelerate the rate of spread. This fuel type is especially susceptible to ignition due to spotting .
Ladder fires consume material between low-level vegetation and tree canopies, such as small trees, downed logs, and vines. Kudzu, Old World climbing fern, and other invasive plants that scale trees may also encourage ladder fires.
Crown, canopy, or aerial fires burn suspended material at the canopy level, such as tall trees, vines, and mosses. The ignition of a crown fire, termed crowning, is dependent on the density of the suspended material, canopy height, canopy continuity, sufficient surface and ladder fires, vegetation moisture content, and weather conditions during the blaze. Stand-replacing fires lit by humans can spread into the Amazon rain forest, damaging ecosystems not particularly suited for heat or arid conditions.
Physical properties
Wildfires occur when all the necessary elements of a fire triangle come together in a susceptible area: an ignition source is brought into contact with a combustible material such as vegetation that is subjected to enough heat and has an adequate supply of oxygen from the ambient air. A high moisture content usually prevents ignition and slows propagation, because higher temperatures are needed to evaporate any water in the material and heat the material to its fire point.
Dense forests usually provide more shade, resulting in lower ambient temperatures and greater humidity, and are therefore less susceptible to wildfires. Less dense material such as grasses and leaves are easier to ignite because they contain less water than denser material such as branches and trunks. Plants continuously lose water by evapotranspiration, but water loss is usually balanced by water absorbed from the soil, humidity, or rain. When this balance is not maintained, often as a consequence of droughts, plants dry out and are therefore more flammable.
A wildfire front is the portion sustaining continuous flaming combustion, where unburned material meets active flames, or the smoldering transition between unburned and burned material. As the front approaches, the fire heats both the surrounding air and woody material through convection and thermal radiation. First, wood is dried as water is vaporized at a temperature of . Next, the pyrolysis of wood at releases flammable gases. Finally, wood can smolder at or, when heated sufficiently, ignite at . Even before the flames of a wildfire arrive at a particular location, heat transfer from the wildfire front warms the air to , which pre-heats and dries flammable materials, causing materials to ignite faster and allowing the fire to spread faster. High-temperature and long-duration surface wildfires may encourage flashover or torching: the drying of tree canopies and their subsequent ignition from below.
Wildfires have a rapid forward rate of spread (FROS) when burning through dense uninterrupted fuels. They can move as fast as in forests and in grasslands. Wildfires can advance tangential to the main front to form a flanking front, or burn in the opposite direction of the main front by backing. They may also spread by jumping or spotting as winds and vertical convection columns carry firebrands (hot wood embers) and other burning materials through the air over roads, rivers, and other barriers that may otherwise act as firebreaks. Torching and fires in tree canopies encourage spotting, and dry ground fuels around a wildfire are especially vulnerable to ignition from firebrands. Spotting can create spot fires as hot embers and firebrands ignite fuels downwind from the fire. In Australian bushfires, spot fires are known to occur as far as from the fire front.
Especially large wildfires may affect air currents in their immediate vicinities by the stack effect: air rises as it is heated, and large wildfires create powerful updrafts that will draw in new, cooler air from surrounding areas in thermal columns. Great vertical differences in temperature and humidity encourage pyrocumulus clouds, strong winds, and fire whirls with the force of tornadoes at speeds of more than . Rapid rates of spread, prolific crowning or spotting, the presence of fire whirls, and strong convection columns signify extreme conditions.
Intensity variations during day and night
Intensity also increases during daytime hours. Burn rates of smoldering logs are up to five times greater during the day due to lower humidity, increased temperatures, and increased wind speeds. Sunlight warms the ground during the day which creates air currents that travel uphill. At night the land cools, creating air currents that travel downhill. Wildfires are fanned by these winds and often follow the air currents over hills and through valleys. Fires in Europe occur frequently during the hours of 12:00 p.m. and 2:00 p.m. Wildfire suppression operations in the United States revolve around a 24-hour fire day that begins at 10:00 a.m. due to the predictable increase in intensity resulting from the daytime warmth.
Climate change effects
Increasing risks due to climate change
Climate change promotes the type of weather that makes wildfires more likely. In some areas, an increase of wildfires has been attributed directly to climate change. Evidence from Earth's past also shows more fire in warmer periods. Climate change increases evapotranspiration. This can cause vegetation and soils to dry out. When a fire starts in an area with very dry vegetation, it can spread rapidly. Higher temperatures can also lengthen the fire season. This is the time of year in which severe wildfires are most likely, particularly in regions where snow is disappearing.
Weather conditions are raising the risks of wildfires. In the United States, about 3 million acres burned each year from 1985 to 1995, but since 2004, between 4 and 10 million acres have burned each year, with the trend continuing to increase. However, the total area burnt by wildfires has decreased worldwide, mostly because savanna has been converted to cropland, so there are fewer trees to burn.
Climate variability including heat waves, droughts, and El Niño, and regional weather patterns, such as high-pressure ridges, can increase the risk and alter the behavior of wildfires dramatically. Years of high precipitation can produce rapid vegetation growth, which when followed by warmer periods can encourage more widespread fires and longer fire seasons. High temperatures dry out the fuel loads and make them more flammable, increasing tree mortality and posing significant risks to global forest health. Since the mid-1980s, in the Western US, earlier snowmelt and associated warming has also been associated with an increase in length and severity of the wildfire season, or the most fire-prone time of the year. A 2019 study indicates that the increase in fire risk in California may be partially attributable to human-induced climate change.
In the summer of 1974–1975 (southern hemisphere), Australia suffered its worst recorded wildfire, when 15% of Australia's land mass suffered "extensive fire damage". Fires that summer burned up an estimated . In Australia, the annual number of hot days (above 35 °C) and very hot days (above 40 °C) has increased significantly in many areas of the country since 1950. The country has always had bushfires but in 2019, the extent and ferocity of these fires increased dramatically. For the first time catastrophic bushfire conditions were declared for Greater Sydney. New South Wales and Queensland declared a state of emergency but fires were also burning in South Australia and Western Australia.
In 2019, extreme heat and dryness caused massive wildfires in Siberia, Alaska, Canary Islands, Australia, and in the Amazon rainforest. The fires in the latter were caused mainly by illegal logging. The smoke from the fires expanded on huge territory including major cities, dramatically reducing air quality.
As of August 2020, the wildfires in that year were 13% worse than in 2019 due primarily to climate change, deforestation and agricultural burning. The Amazon rainforest's existence is threatened by fires. Record-breaking wildfires in 2021 occurred in Turkey, Greece and Russia, thought to be linked to climate change.
Carbon dioxide and other emissions from fires
The carbon released from wildfires can add to greenhouse gas concentrations. Climate models do not yet fully reflect this feedback.
Wildfires release large amounts of carbon dioxide, black and brown carbon particles, and ozone precursors such as volatile organic compounds and nitrogen oxides (NOx) into the atmosphere. These emissions affect radiation, clouds, and climate on regional and even global scales. Wildfires also emit substantial amounts of semi-volatile organic species that can partition from the gas phase to form secondary organic aerosol (SOA) over hours to days after emission. In addition, the formation of the other pollutants as the air is transported can lead to harmful exposures for populations in regions far away from the wildfires. While direct emissions of harmful pollutants can affect first responders and residents, wildfire smoke can also be transported over long distances and impact air quality across local, regional, and global scales.The health effects of wildfire smoke, such as worsening cardiovascular and respiratory conditions, extend beyond immediate exposure, contributing to nearly 16,000 annual deaths, a number expected to rise to 30,000 by 2050. The economic impact is also significant, with projected costs reaching $240 billion annually by 2050, surpassing other climate-related damages.
Over the past century, wildfires have accounted for 20–25% of global carbon emissions, the remainder from human activities. Global carbon emissions from wildfires through August 2020 equaled the average annual emissions of the European Union. In 2020, the carbon released by California's wildfires was significantly larger than the state's other carbon emissions.
Forest fires in Indonesia in 1997 were estimated to have released between 0.81 and 2.57 gigatonnes (0.89 and 2.83 billion short tons) of CO2 into the atmosphere, which is between 13–40% of the annual global carbon dioxide emissions from burning fossil fuels.
In June and July 2019, fires in the Arctic emitted more than 140 megatons of carbon dioxide, according to an analysis by CAMS. To put that into perspective this amounts to the same amount of carbon emitted by 36 million cars in a year. The recent wildfires and their massive CO2 emissions mean that it will be important to take them into consideration when implementing measures for reaching greenhouse gas reduction targets accorded with the Paris climate agreement. Due to the complex oxidative chemistry occurring during the transport of wildfire smoke in the atmosphere, the toxicity of emissions was indicated to increase over time.
Atmospheric models suggest that these concentrations of sooty particles could increase absorption of incoming solar radiation during winter months by as much as 15%. The Amazon is estimated to hold around 90 billion tons of carbon. As of 2019, the earth's atmosphere has 415 parts per million of carbon, and the destruction of the Amazon would add about 38 parts per million.
Some research has shown wildfire smoke can have a cooling effect.
Research in 2007 stated that black carbon in snow changed temperature three times more than atmospheric carbon dioxide. As much as 94 percent of Arctic warming may be caused by dark carbon on snow that initiates melting. The dark carbon comes from fossil fuels burning, wood and other biofuels, and forest fires. Melting can occur even at low concentrations of dark carbon (below five parts per billion).
Prevention
Wildfire prevention refers to the preemptive methods aimed at reducing the risk of fires as well as lessening its severity and spread. Prevention techniques aim to manage air quality, maintain ecological balances, protect resources, and to affect future fires. Prevention policies must consider the role that humans play in wildfires, since, for example, 95% of forest fires in Europe are related to human involvement.
Wildfire prevention programs around the world may employ techniques such as wildland fire use (WFU) and prescribed or controlled burns. Wildland fire use refers to any fire of natural causes that is monitored but allowed to burn. Controlled burns are fires ignited by government agencies under less dangerous weather conditions. Other objectives can include maintenance of healthy forests, rangelands, and wetlands, and support of ecosystem diversity.
Strategies for wildfire prevention, detection, control and suppression have varied over the years. One common and inexpensive technique to reduce the risk of uncontrolled wildfires is controlled burning: intentionally igniting smaller less-intense fires to minimize the amount of flammable material available for a potential wildfire. Vegetation may be burned periodically to limit the accumulation of plants and other debris that may serve as fuel, while also maintaining high species diversity. While other people claim that controlled burns and a policy of allowing some wildfires to burn is the cheapest method and an ecologically appropriate policy for many forests, they tend not to take into account the economic value of resources that are consumed by the fire, especially merchantable timber. Some studies conclude that while fuels may also be removed by logging, such thinning treatments may not be effective at reducing fire severity under extreme weather conditions.
Building codes in fire-prone areas typically require that structures be built of flame-resistant materials and a defensible space be maintained by clearing flammable materials within a prescribed distance from the structure. Communities in the Philippines also maintain fire lines wide between the forest and their village, and patrol these lines during summer months or seasons of dry weather. Continued residential development in fire-prone areas and rebuilding structures destroyed by fires has been met with criticism. The ecological benefits of fire are often overridden by the economic and safety benefits of protecting structures and human life.
Detection
The demand for timely, high-quality fire information has increased in recent years. Fast and effective detection is a key factor in wildfire fighting. Early detection efforts were focused on early response, accurate results in both daytime and nighttime, and the ability to prioritize fire danger. Fire lookout towers were used in the United States in the early 20th century and fires were reported using telephones, carrier pigeons, and heliographs. Aerial and land photography using instant cameras were used in the 1950s until infrared scanning was developed for fire detection in the 1960s. However, information analysis and delivery was often delayed by limitations in communication technology. Early satellite-derived fire analyses were hand-drawn on maps at a remote site and sent via overnight mail to the fire manager. During the Yellowstone fires of 1988, a data station was established in West Yellowstone, permitting the delivery of satellite-based fire information in approximately four hours.
Public hotlines, fire lookouts in towers, and ground and aerial patrols can be used as a means of early detection of forest fires. However, accurate human observation may be limited by operator fatigue, time of day, time of year, and geographic location. Electronic systems have gained popularity in recent years as a possible resolution to human operator error. These systems may be semi- or fully automated and employ systems based on the risk area and degree of human presence, as suggested by GIS data analyses. An integrated approach of multiple systems can be used to merge satellite data, aerial imagery, and personnel position via Global Positioning System (GPS) into a collective whole for near-realtime use by wireless Incident Command Centers.
Local sensor networks
A small, high risk area that features thick vegetation, a strong human presence, or is close to a critical urban area can be monitored using a local sensor network. Detection systems may include wireless sensor networks that act as automated weather systems: detecting temperature, humidity, and smoke. These may be battery-powered, solar-powered, or tree-rechargeable: able to recharge their battery systems using the small electrical currents in plant material. Larger, medium-risk areas can be monitored by scanning towers that incorporate fixed cameras and sensors to detect smoke or additional factors such as the infrared signature of carbon dioxide produced by fires. Additional capabilities such as night vision, brightness detection, and color change detection may also be incorporated into sensor arrays.
The Department of Natural Resources signed a contract with PanoAI for the installation of 360 degree 'rapid detection' cameras around the Pacific northwest, which are mounted on cell towers and are capable of 24/7 monitoring of a 15 mile radius. Additionally, Sensaio Tech, based in Brazil and Toronto, has released a sensor device that continuously monitors 14 different variables common in forests, ranging from soil temperature to salinity. This information is connected live back to clients through dashboard visualizations, while mobile notifications are provided regarding dangerous levels.
Satellite and aerial monitoring
Satellite and aerial monitoring through the use of planes, helicopter, or UAVs can provide a wider view and may be sufficient to monitor very large, low risk areas. These more sophisticated systems employ GPS and aircraft-mounted infrared or high-resolution visible cameras to identify and target wildfires. Satellite-mounted sensors such as Envisat's Advanced Along Track Scanning Radiometer and European Remote-Sensing Satellite's Along-Track Scanning Radiometer can measure infrared radiation emitted by fires, identifying hot spots greater than . The National Oceanic and Atmospheric Administration's Hazard Mapping System combines remote-sensing data from satellite sources such as Geostationary Operational Environmental Satellite (GOES), Moderate-Resolution Imaging Spectroradiometer (MODIS), and Advanced Very High Resolution Radiometer (AVHRR) for detection of fire and smoke plume locations. However, satellite detection is prone to offset errors, anywhere from for MODIS and AVHRR data and up to for GOES data. Satellites in geostationary orbits may become disabled, and satellites in polar orbits are often limited by their short window of observation time. Cloud cover and image resolution may also limit the effectiveness of satellite imagery. Global Forest Watch provides detailed daily updates on fire alerts.
In 2015 a new fire detection tool is in operation at the U.S. Department of Agriculture (USDA) Forest Service (USFS) which uses data from the Suomi National Polar-orbiting Partnership (NPP) satellite to detect smaller fires in more detail than previous space-based products. The high-resolution data is used with a computer model to predict how a fire will change direction based on weather and land conditions.
In 2014, an international campaign was organized in South Africa's Kruger National Park to validate fire detection products including the new VIIRS active fire data. In advance of that campaign, the Meraka Institute of the Council for Scientific and Industrial Research in Pretoria, South Africa, an early adopter of the VIIRS 375 m fire product, put it to use during several large wildfires in Kruger.
Since 2021 NASA has provided active fire locations in near real-time via the Fire Information for Resource Management System (FIRMS).
Between 2022–2023, wildfires throughout North America prompted an uptake in the delivery and design of various technologies using artificial intelligence for early detection, prevention, and prediction of wildfires.
Suppression
Wildfire suppression depends on the technologies available in the area in which the wildfire occurs. In less developed nations the techniques used can be as simple as throwing sand or beating the fire with sticks or palm fronds. In more advanced nations, the suppression methods vary due to increased technological capacity. Silver iodide can be used to encourage snow fall, while fire retardants and water can be dropped onto fires by unmanned aerial vehicles, planes, and helicopters. Complete fire suppression is no longer an expectation, but the majority of wildfires are often extinguished before they grow out of control. While more than 99% of the 10,000 new wildfires each year are contained, escaped wildfires under extreme weather conditions are difficult to suppress without a change in the weather. Wildfires in Canada and the US burn an average of per year.
Above all, fighting wildfires can become deadly. A wildfire's burning front may also change direction unexpectedly and jump across fire breaks. Intense heat and smoke can lead to disorientation and loss of appreciation of the direction of the fire, which can make fires particularly dangerous. For example, during the 1949 Mann Gulch fire in Montana, United States, thirteen smokejumpers died when they lost their communication links, became disoriented, and were overtaken by the fire. In the Australian February 2009 Victorian bushfires, at least 173 people died and over 2,029 homes and 3,500 structures were lost when they became engulfed by wildfire.
Costs of wildfire suppression
The suppression of wild fires takes up a large amount of a country's gross domestic product which directly affects the country's economy. While costs vary wildly from year to year, depending on the severity of each fire season, in the United States, local, state, federal and tribal agencies collectively spend tens of billions of dollars annually to suppress wildfires. In the United States, it was reported that approximately $6 billion was spent between 2004–2008 to suppress wildfires in the country. In California, the U.S. Forest Service spends about $200 million per year to suppress 98% of wildfires and up to $1 billion to suppress the other 2% of fires that escape initial attack and become large.
Wildland firefighting safety
Wildland fire fighters face several life-threatening hazards including heat stress, fatigue, smoke and dust, as well as the risk of other injuries such as burns, cuts and scrapes, animal bites, and even rhabdomyolysis. Between 2000 and 2016, more than 350 wildland firefighters died on-duty.
Especially in hot weather conditions, fires present the risk of heat stress, which can entail feeling heat, fatigue, weakness, vertigo, headache, or nausea. Heat stress can progress into heat strain, which entails physiological changes such as increased heart rate and core body temperature. This can lead to heat-related illnesses, such as heat rash, cramps, exhaustion or heat stroke. Various factors can contribute to the risks posed by heat stress, including strenuous work, personal risk factors such as age and fitness, dehydration, sleep deprivation, and burdensome personal protective equipment. Rest, cool water, and occasional breaks are crucial to mitigating the effects of heat stress.
Smoke, ash, and debris can also pose serious respiratory hazards for wildland firefighters. The smoke and dust from wildfires can contain gases such as carbon monoxide, sulfur dioxide and formaldehyde, as well as particulates such as ash and silica. To reduce smoke exposure, wildfire fighting crews should, whenever possible, rotate firefighters through areas of heavy smoke, avoid downwind firefighting, use equipment rather than people in holding areas, and minimize mop-up. Camps and command posts should also be located upwind of wildfires. Protective clothing and equipment can also help minimize exposure to smoke and ash.
Firefighters are also at risk of cardiac events including strokes and heart attacks. Firefighters should maintain good physical fitness. Fitness programs, medical screening and examination programs which include stress tests can minimize the risks of firefighting cardiac problems. Other injury hazards wildland firefighters face include slips, trips, falls, burns, scrapes, and cuts from tools and equipment, being struck by trees, vehicles, or other objects, plant hazards such as thorns and poison ivy, snake and animal bites, vehicle crashes, electrocution from power lines or lightning storms, and unstable building structures.
Fire retardants
Fire retardants are used to slow wildfires by inhibiting combustion. They are aqueous solutions of ammonium phosphates and ammonium sulfates, as well as thickening agents. The decision to apply retardant depends on the magnitude, location and intensity of the wildfire. In certain instances, fire retardant may also be applied as a precautionary fire defense measure.
Typical fire retardants contain the same agents as fertilizers. Fire retardants may also affect water quality through leaching, eutrophication, or misapplication. Fire retardant's effects on drinking water remain inconclusive. Dilution factors, including water body size, rainfall, and water flow rates lessen the concentration and potency of fire retardant. Wildfire debris (ash and sediment) clog rivers and reservoirs increasing the risk for floods and erosion that ultimately slow and/or damage water treatment systems. There is continued concern of fire retardant effects on land, water, wildlife habitats, and watershed quality, additional research is needed. However, on the positive side, fire retardant (specifically its nitrogen and phosphorus components) has been shown to have a fertilizing effect on nutrient-deprived soils and thus creates a temporary increase in vegetation.
Modeling
Impacts on the natural environment
On the atmosphere
Most of Earth's weather and air pollution resides in the troposphere, the part of the atmosphere that extends from the surface of the planet to a height of about . The vertical lift of a severe thunderstorm or pyrocumulonimbus can be enhanced in the area of a large wildfire, which can propel smoke, soot (black carbon), and other particulate matter as high as the lower stratosphere. Previously, prevailing scientific theory held that most particles in the stratosphere came from volcanoes, but smoke and other wildfire emissions have been detected from the lower stratosphere. Pyrocumulus clouds can reach over wildfires. Satellite observation of smoke plumes from wildfires revealed that the plumes could be traced intact for distances exceeding . Computer-aided models such as CALPUFF may help predict the size and direction of wildfire-generated smoke plumes by using atmospheric dispersion modeling.
Wildfires can affect local atmospheric pollution, and release carbon in the form of carbon dioxide. Wildfire emissions contain fine particulate matter which can cause cardiovascular and respiratory problems. Increased fire byproducts in the troposphere can increase ozone concentrations beyond safe levels.
On ecosystems
Wildfires are common in climates that are sufficiently moist to allow the growth of vegetation but feature extended dry, hot periods. Such places include the vegetated areas of Australia and Southeast Asia, the veld in southern Africa, the fynbos in the Western Cape of South Africa, the forested areas of the United States and Canada, and the Mediterranean Basin.
High-severity wildfire creates complex early seral forest habitat (also called "snag forest habitat"), which often has higher species richness and diversity than unburned old forest. Plant and animal species in most types of North American forests evolved with fire, and many of these species depend on wildfires, and particularly high-severity fires, to reproduce and grow. Fire helps to return nutrients from plant matter back to the soil. The heat from fire is necessary to the germination of certain types of seeds, and the snags (dead trees) and early successional forests created by high-severity fire create habitat conditions that are beneficial to wildlife. Early successional forests created by high-severity fire support some of the highest levels of native biodiversity found in temperate conifer forests. Post-fire logging has no ecological benefits and many negative impacts; the same is often true for post-fire seeding. The exclusion of wildfires can contribute to vegetation regime shifts, such as woody plant encroachment.
Although some ecosystems rely on naturally occurring fires to regulate growth, some ecosystems suffer from too much fire, such as the chaparral in southern California and lower-elevation deserts in the American Southwest. The increased fire frequency in these ordinarily fire-dependent areas has upset natural cycles, damaged native plant communities, and encouraged the growth of non-native weeds. Invasive species, such as Lygodium microphyllum and Bromus tectorum, can grow rapidly in areas that were damaged by fires. Because they are highly flammable, they can increase the future risk of fire, creating a positive feedback loop that increases fire frequency and further alters native vegetation communities.
In the Amazon rainforest, drought, logging, cattle ranching practices, and slash-and-burn agriculture damage fire-resistant forests and promote the growth of flammable brush, creating a cycle that encourages more burning. Fires in the rainforest threaten its collection of diverse species and produce large amounts of CO2. Also, fires in the rainforest, along with drought and human involvement, could damage or destroy more than half of the Amazon rainforest by 2030. Wildfires generate ash, reduce the availability of organic nutrients, and cause an increase in water runoff, eroding other nutrients and creating flash flood conditions. A 2003 wildfire in the North Yorkshire Moors burned off of heather and the underlying peat layers. Afterwards, wind erosion stripped the ash and the exposed soil, revealing archaeological remains dating to 10,000 BC. Wildfires can also have an effect on climate change, increasing the amount of carbon released into the atmosphere and inhibiting vegetation growth, which affects overall carbon uptake by plants.
On waterways
Debris and chemical runoff into waterways after wildfires can make drinking water sources unsafe. Though it is challenging to quantify the impacts of wildfires on surface water quality, research suggests that the concentration of many pollutants increases post-fire. The impacts occur during active burning and up to years later. Increases in nutrients and total suspended sediments can happen within a year while heavy metal concentrations may peak 1–2 years after a wildfire.
Benzene is one of many chemicals that have been found in drinking water systems after wildfires. Benzene can permeate certain plastic pipes and thus require long times to be removed from the water distribution infrastructure. Researchers estimated that, in worst case scenarios, more than 286 days of constant flushing of a contaminated HDPE service line were needed to reduce benzene below safe drinking water limits. Temperature increases caused by fires, including wildfires, can cause plastic water pipes to generate toxic chemicals such as benzene.
On plant and animals
Impacts on humans
Wildfire risk is the chance that a wildfire will start in or reach a particular area and the potential loss of human values if it does. Risk is dependent on variable factors such as human activities, weather patterns, availability of wildfire fuels, and the availability or lack of resources to suppress a fire. Wildfires have continually been a threat to human populations. However, human-induced geographic and climatic changes are exposing populations more frequently to wildfires and increasing wildfire risk. It is speculated that the increase in wildfires arises from a century of wildfire suppression coupled with the rapid expansion of human developments into fire-prone wildlands. Wildfires are naturally occurring events that aid in promoting forest health. Global warming and climate changes are causing an increase in temperatures and more droughts nationwide which contributes to an increase in wildfire risk.
Airborne hazards
The most noticeable adverse effect of wildfires is the destruction of property. However, hazardous chemicals released also significantly impact human health.
Wildfire smoke is composed primarily of carbon dioxide and water vapor. Other common components present in lower concentrations are carbon monoxide, formaldehyde, acrolein, polyaromatic hydrocarbons, and benzene. Small airborne particulates (in solid form or liquid droplets) are also present in smoke and ash debris. 80–90% of wildfire smoke, by mass, is within the fine particle size class of 2.5 micrometers in diameter or smaller.
Carbon dioxide in smoke poses a low health risk due to its low toxicity. Rather, carbon monoxide and fine particulate matter, particularly 2.5 μm in diameter and smaller, have been identified as the major health threats. High levels of heavy metals, including lead, arsenic, cadmium, and copper were found in the ash debris following the 2007 Californian wildfires. A national clean-up campaign was organised in fear of the health effects from exposure. In the devastating California Camp Fire (2018) that killed 85 people, lead levels increased by around 50 times in the hours following the fire at a site nearby (Chico). Zinc concentration also increased significantly in Modesto, 150 miles away. Heavy metals such as manganese and calcium were found in numerous California fires as well. Other chemicals are considered to be significant hazards but are found in concentrations that are too low to cause detectable health effects.
The degree of wildfire smoke exposure to an individual is dependent on the length, severity, duration, and proximity of the fire. People are exposed directly to smoke via the respiratory tract through inhalation of air pollutants. Indirectly, communities are exposed to wildfire debris that can contaminate soil and water supplies.
The U.S. Environmental Protection Agency (EPA) developed the air quality index (AQI), a public resource that provides national air quality standard concentrations for common air pollutants. The public can use it to determine their exposure to hazardous air pollutants based on visibility range.
Health effects
Wildfire smoke contains particulates that may have adverse effects upon the human respiratory system. Evidence of the health effects should be relayed to the public so that exposure may be limited. The evidence can also be used to influence policy to promote positive health outcomes.
Inhalation of smoke from a wildfire can be a health hazard. Wildfire smoke is composed of combustion products i.e. carbon dioxide, carbon monoxide, water vapor, particulate matter, organic chemicals, nitrogen oxides and other compounds. The principal health concern is the inhalation of particulate matter and carbon monoxide.
Particulate matter (PM) is a type of air pollution made up of particles of dust and liquid droplets. They are characterized into three categories based on particle diameter: coarse PM, fine PM, and ultrafine PM. Coarse particles are between 2.5 micrometers and 10 micrometers, fine particles measure 0.1 to 2.5 micrometers, and ultrafine particle are less than 0.1 micrometer. lmpact on the body upon inhalation varies by size. Coarse PM is filtered by the upper airways and can accumulate and cause pulmonary inflammation. This can result in eye and sinus irritation as well as sore throat and coughing. Coarse PM is often composed of heavier and more toxic materials that lead to short-term effects with stronger impact.
Smaller PM moves further into the respiratory system creating issues deep into the lungs and the bloodstream. In asthma patients, PM2.5 causes inflammation but also increases oxidative stress in the epithelial cells. These particulates also cause apoptosis and autophagy in lung epithelial cells. Both processes damage the cells and impact cell function. This damage impacts those with respiratory conditions such as asthma where the lung tissues and function are already compromised. Particulates less than 0.1 micrometer are called ultrafine particle (UFP). It is a major component of wildfire smoke. UFP can enter the bloodstream like PM2.5–0.1 however studies show that it works into the blood much quicker. The inflammation and epithelial damage done by UFP has also shown to be much more severe. PM2.5 is of the largest concern in regards to wildfire. This is particularly hazardous to the very young, elderly and those with chronic conditions such as asthma, chronic obstructive pulmonary disease (COPD), cystic fibrosis and cardiovascular conditions. The illnesses most commonly associated with exposure to fine PM from wildfire smoke are bronchitis, exacerbation of asthma or COPD, and pneumonia. Symptoms of these complications include wheezing and shortness of breath and cardiovascular symptoms include chest pain, rapid heart rate and fatigue.
Asthma exacerbation
Several epidemiological studies have demonstrated a close association between air pollution and respiratory allergic diseases such as bronchial asthma.
An observational study of smoke exposure related to the 2007 San Diego wildfires revealed an increase both in healthcare utilization and respiratory diagnoses, especially asthma among the group sampled. Projected climate scenarios of wildfire occurrences predict significant increases in respiratory conditions among young children. PM triggers a series of biological processes including inflammatory immune response, oxidative stress, which are associated with harmful changes in allergic respiratory diseases.
Although some studies demonstrated no significant acute changes in lung function among people with asthma related to PM from wildfires, a possible explanation for these counterintuitive findings is the increased use of quick-relief medications, such as inhalers, in response to elevated levels of smoke among those already diagnosed with asthma.
There is consistent evidence between wildfire smoke and the exacerbation of asthma.
Asthma is one of the most common chronic disease among children in the United States, affecting an estimated 6.2 million children. Research on asthma risk focuses specifically on the risk of air pollution during the gestational period. Several pathophysiology processes are involved in this. Considerable airway development occurs during the 2nd and 3rd trimesters and continues until 3 years of age. It is hypothesized that exposure to these toxins during this period could have consequential effects, as the epithelium of the lungs during this time could have increased permeability to toxins. Exposure to air pollution during parental and pre-natal stage could induce epigenetic changes which are responsible for the development of asthma. Studies have found significant association between PM2.5, NO2 and development of asthma during childhood despite heterogeneity among studies. Furthermore, maternal exposure to chronic stressors is most likely present in distressed communities, and as this can be correlated with childhood asthma, it may further explain links between early childhood exposure to air pollution, neighborhood poverty, and childhood risk.
Carbon monoxide danger
Carbon monoxide (CO) is a colorless, odorless gas that can be found at the highest concentration at close proximity to a smoldering fire. Thus, it is a serious threat to the health of wildfire firefighters. CO in smoke can be inhaled into the lungs where it is absorbed into the bloodstream and reduces oxygen delivery to the body's vital organs. At high concentrations, it can cause headaches, weakness, dizziness, confusion, nausea, disorientation, visual impairment, coma, and even death. Even at lower concentrations, such as those found at wildfires, individuals with cardiovascular disease may experience chest pain and cardiac arrhythmia. A recent study tracking the number and cause of wildfire firefighter deaths from 1990 to 2006 found that 21.9% of the deaths occurred from heart attacks.
Another important and somewhat less obvious health effect of wildfires is psychiatric diseases and disorders. Both adults and children from various countries who were directly and indirectly affected by wildfires were found to demonstrate different mental conditions linked to their experience with the wildfires. These include post-traumatic stress disorder (PTSD), depression, anxiety, and phobias.
Epidemiology
The Western US has seen an increase in both the frequency and intensity of wildfires over the last several decades. This has been attributed to the arid climate of there and the effects of global warming. An estimated 46 million people were exposed to wildfire smoke from 2004 to 2009 in the Western US. Evidence has demonstrated that wildfire smoke can increase levels of airborne particulate.
The EPA has defined acceptable concentrations of PM in the air, through the National Ambient Air Quality Standards and monitoring of ambient air quality has been mandated. Due to these monitoring programs and the incidence of several large wildfires near populated areas, epidemiological studies have been conducted and demonstrate an association between human health effects and an increase in fine particulate matter due to wildfire smoke.
An increase in PM smoke emitted from the Hayman fire in Colorado in June 2002, was associated with an increase in respiratory symptoms in patients with COPD. Looking at the wildfires in Southern California in 2003, investigators have shown an increase in hospital admissions due to asthma symptoms while being exposed to peak concentrations of PM in smoke. Another epidemiological study found a 7.2% (95% confidence interval: 0.25%, 15%) increase in risk of respiratory related hospital admissions during smoke wave days with high wildfire-specific particulate matter 2.5 compared to matched non-smoke-wave days.
Children participating in the Children's Health Study were also found to have an increase in eye and respiratory symptoms, medication use and physician visits. Mothers who were pregnant during the fires gave birth to babies with a slightly reduced average birth weight compared to those who were not exposed. Suggesting that pregnant women may also be at greater risk to adverse effects from wildfire. Worldwide, it is estimated that 339,000 people die due to the effects of wildfire smoke each year.
Besides the size of PM, their chemical composition should also be considered. Antecedent studies have demonstrated that the chemical composition of PM2.5 from wildfire smoke can yield different estimates of human health outcomes as compared to other sources of smoke such as solid fuels.
Post-fire risks
After a wildfire, hazards remain. Residents returning to their homes may be at risk from falling fire-weakened trees. Humans and pets may also be harmed by falling into ash pits. The Intergovernmental Panel on Climate Change (IPCC) also reports that wildfires cause significant damage to electric systems, especially in dry regions.
Chemically contaminated drinking water, at levels of hazardous waste concern, is a growing problem. In particular, hazardous waste scale chemical contamination of buried water systems was first discovered in the U.S. in 2017, and has since been increasingly documented in Hawaii, Colorado, and Oregon after wildfires. In 2021, Canadian authorities adapted their post-fire public safety investigation approaches in British Columbia to screen for this risk, but have not found it as of 2023. Another challenge is that private drinking wells and the plumbing within a building can also become chemically contaminated and unsafe. Households experience a wide-variety of significant economic and health impacts related to this contaminated water. Evidence-based guidance on how to inspect and test wildfire impacted wells and building water systems was developed for the first time in 2020. In Paradise, California, for example, the 2018 Camp Fire caused more than $150 million dollars worth of damage. This required almost a year of time to decontaminate and repair the municipal drinking water system from wildfire damage.
The source of this contamination was first proposed after the 2018 Camp Fire in California as originating from thermally degraded plastics in water systems, smoke and vapors entering depressurized plumbing, and contaminated water in buildings being sucked into the municipal water system. In 2020, it was first shown that thermal degradation of plastic drinking water materials was one potential contamination source. In 2023, the second theory was confirmed where contamination could be sucked into pipes that lost water pressure.
Other post-fire risks, can increase if other extreme weather follows. For example, wildfires make soil less able to absorb precipitation, so heavy rainfall can result in more severe flooding and damages like mud slides.
At-risk groups
Firefighters
Firefighters are at greatest risk for acute and chronic health effects resulting from wildfire smoke exposure. Some of the most common health conditions that firefighters acquire from prolonged smoke inhalation include cardiovascular and respiratory diseases. For example, wildland firefighters can get hypoxia, which is a condition in which the body does not receive enough oxygen. Due to firefighters' occupational duties, they are frequently exposed to hazardous chemicals at close proximity for longer periods of time. A case study on the exposure of wildfire smoke among wildland firefighters shows that firefighters are exposed to significant levels of carbon monoxide and respiratory irritants above OSHA-permissible exposure limits (PEL) and ACGIH threshold limit values (TLV). 5–10% are overexposed.
Between 2001 and 2012, over 200 fatalities occurred among wildland firefighters. In addition to heat and chemical hazards, firefighters are also at risk for electrocution from power lines; injuries from equipment; slips, trips, and falls; injuries from vehicle rollovers; heat-related illness; insect bites and stings; stress; and rhabdomyolysis.
Residents
Residents in communities surrounding wildfires are exposed to lower concentrations of chemicals, but they are at a greater risk for indirect exposure through water or soil contamination. Exposure to residents is greatly dependent on individual susceptibility. Vulnerable persons such as children (ages 0–4), the elderly (ages 65 and older), smokers, and pregnant women are at an increased risk due to their already compromised body systems, even when the exposures are present at low chemical concentrations and for relatively short exposure periods. They are also at risk for future wildfires and may move away to areas they consider less risky.
Wildfires affect large numbers of people in Western Canada and the United States. In California alone, more than 350,000 people live in towns and cities in "very high fire hazard severity zones".
Direct risks to building residents in fire-prone areas can be moderated through design choices such as choosing fire-resistant vegetation, maintaining landscaping to avoid debris accumulation and to create firebreaks, and by selecting fire-retardant roofing materials. Potential compounding issues with poor air quality and heat during warmer months may be addressed with MERV 11 or higher outdoor air filtration in building ventilation systems, mechanical cooling, and a provision of a refuge area with additional air cleaning and cooling, if needed.
History
The first evidence of wildfires is fossils of the giant fungi Prototaxites preserved as charcoal, discovered in South Wales and Poland, dating to the Silurian period (about ). Smoldering surface fires started to occur sometime before the Early Devonian period . Low atmospheric oxygen during the Middle and Late Devonian was accompanied by a decrease in charcoal abundance. Additional charcoal evidence suggests that fires continued through the Carboniferous period. Later, the overall increase of atmospheric oxygen from 13% in the Late Devonian to 30–31% by the Late Permian was accompanied by a more widespread distribution of wildfires. Later, a decrease in wildfire-related charcoal deposits from the late Permian to the Triassic periods is explained by a decrease in oxygen levels.
Wildfires during the Paleozoic and Mesozoic periods followed patterns similar to fires that occur in modern times. Surface fires driven by dry seasons are evident in Devonian and Carboniferous progymnosperm forests. Lepidodendron forests dating to the Carboniferous period have charred peaks, evidence of crown fires. In Jurassic gymnosperm forests, there is evidence of high frequency, light surface fires. The increase of fire activity in the late Tertiary is possibly due to the increase of C4-type grasses. As these grasses shifted to more mesic habitats, their high flammability increased fire frequency, promoting grasslands over woodlands. However, fire-prone habitats may have contributed to the prominence of trees such as those of the genera Eucalyptus, Pinus and Sequoia, which have thick bark to withstand fires and employ pyriscence.
Human involvement
The human use of fire for agricultural and hunting purposes during the Paleolithic and Mesolithic ages altered pre-existing landscapes and fire regimes. Woodlands were gradually replaced by smaller vegetation that facilitated travel, hunting, seed-gathering and planting. In recorded human history, minor allusions to wildfires were mentioned in the Bible and by classical writers such as Homer. However, while ancient Hebrew, Greek, and Roman writers were aware of fires, they were not very interested in the uncultivated lands where wildfires occurred. Wildfires were used in battles throughout human history as early thermal weapons. From the Middle Ages, accounts were written of occupational burning as well as customs and laws that governed the use of fire. In Germany, regular burning was documented in 1290 in the Odenwald and in 1344 in the Black Forest. In the 14th century Sardinia, firebreaks were used for wildfire protection. In Spain during the 1550s, sheep husbandry was discouraged in certain provinces by Philip II due to the harmful effects of fires used in transhumance. As early as the 17th century, Native Americans were observed using fire for many purposes including cultivation, signaling, and warfare. Scottish botanist David Douglas noted the native use of fire for tobacco cultivation, to encourage deer into smaller areas for hunting purposes, and to improve foraging for honey and grasshoppers. Charcoal found in sedimentary deposits off the Pacific coast of Central America suggests that more burning occurred in the 50 years before the Spanish colonization of the Americas than after the colonization. In the post-World War II Baltic region, socio-economic changes led more stringent air quality standards and bans on fires that eliminated traditional burning practices. In the mid-19th century, explorers from observed Australian Aborigines using fire for ground clearing, hunting, and regeneration of plant food in a method later named fire-stick farming. Such careful use of fire has been employed for centuries in lands protected by Kakadu National Park to encourage biodiversity.
Wildfires typically occur during periods of increased temperature and drought. An increase in fire-related debris flow in alluvial fans of northeastern Yellowstone National Park was linked to the period between AD 1050 and 1200, coinciding with the Medieval Warm Period. However, human influence caused an increase in fire frequency. Dendrochronological fire scar data and charcoal layer data in Finland suggests that, while many fires occurred during severe drought conditions, an increase in the number of fires during 850 BC and 1660 AD can be attributed to human influence. Charcoal evidence from the Americas suggested a general decrease in wildfires between 1 AD and 1750 compared to previous years. However, a period of increased fire frequency between 1750 and 1870 was suggested by charcoal data from North America and Asia, attributed to human population growth and influences such as land clearing practices. This period was followed by an overall decrease in burning in the 20th century, linked to the expansion of agriculture, increased livestock grazing, and fire prevention efforts. A meta-analysis found that 17 times more land burned annually in California before 1800 compared to recent decades (1,800,000 hectares/year compared to 102,000 hectares/year).
According to a paper published in the journal Science, the number of natural and human-caused fires decreased by 24.3% between 1998 and 2015. Researchers explain this as a transition from nomadism to settled lifestyle and intensification of agriculture that lead to a drop in the use of fire for land clearing.
Increases of certain tree species (i.e. conifers) over others (i.e. deciduous trees) can increase wildfire risk, especially if these trees are also planted in monocultures.
Some invasive species, moved in by humans (i.e., for the pulp and paper industry) have in some cases also increased the intensity of wildfires. Examples include species such as Eucalyptus in California and gamba grass in Australia.
Society and culture
Wildfires have a place in many cultures. "To spread like wildfire" is a common idiom in English, meaning something that "quickly affects or becomes known by more and more people".
Wildfire activity has been attributed as a major factor in the development of Ancient Greece. In modern Greece, as in many other regions, it is the most common disaster caused by a natural hazard and figures prominently in the social and economic lives of its people.
In 1937, U.S. President Franklin D. Roosevelt initiated a nationwide fire prevention campaign, highlighting the role of human carelessness in forest fires. Later posters of the program featured Uncle Sam, characters from the Disney movie Bambi, and the official mascot of the U.S. Forest Service, Smokey Bear. The Smokey Bear fire prevention campaign has yielded one of the most popular characters in the United States; for many years there was a living Smokey Bear mascot, and it has been commemorated on postage stamps.
There are also significant indirect or second-order societal impacts from wildfire, such as demands on utilities to prevent power transmission equipment from becoming ignition sources, and the cancelation or nonrenewal of homeowners insurance for residents living in wildfire-prone areas.
| Physical sciences | Natural disasters | null |
56109 | https://en.wikipedia.org/wiki/Brown%20rat | Brown rat | The brown rat (Rattus norvegicus), also known as the common rat, street rat, sewer rat, wharf rat, Hanover rat, Norway rat and Norwegian rat, is a widespread species of common rat. One of the largest muroids, it is a brown or grey rodent with a body length of up to long, and a tail slightly shorter than that. It weighs between . Thought to have originated in northern China and neighbouring areas, this rodent has now spread to all continents except Antarctica, and is the dominant rat in Europe and much of North America. With rare exceptions, the brown rat lives wherever humans live, particularly in urban areas.
Selective breeding of the brown rat has produced the fancy rat (rats kept as pets), as well as the laboratory rat (rats used as model organisms in biological research). Both fancy rats and laboratory rats are of the domesticated subspecies Rattus norvegicus domestica. Studies of wild rats in New York City have shown that populations living in different neighborhoods can evolve distinct genomic profiles over time, by slowly accruing different traits.
Naming and etymology
The brown rat was originally called the "Hanover rat" by people wishing to link problems in 18th-century England with the House of Hanover. It is not known for certain why the brown rat is named Rattus norvegicus (Norwegian rat), as it did not originate from Norway. However, the English naturalist John Berkenhout, author of the 1769 book Outlines of the Natural History of Great Britain, is most likely responsible for popularizing the misnomer. Berkenhout gave the brown rat the binomial name Rattus norvegicus, believing it had migrated to England from Norwegian ships in 1728.
By the early to the middle part of the 19th century, British academics believed that the brown rat was not native to Norway, hypothesizing (incorrectly) that it may have come from Ireland, Gibraltar or across the English Channel with William the Conqueror. As early as 1850, however, a new hypothesis of the rat's origins was beginning to develop. The British novelist Charles Dickens acknowledged this in his weekly journal, All the Year Round, writing:
It is frequently called, in books and otherwise, the 'Norway rat', and it is said to have been imported into this country in a ship-load of timber from Norway. Against this hypothesis stands the fact that when the brown rat had become common in this country, it was unknown in Norway, although there was a small animal like a rat, but really a lemming, which made its home there.
Academics began to prefer this etymology of the brown rat towards the end of the 19th century, as seen in the 1895 text Natural History by American scholar Alfred Henry Miles:
The brown rat is the species common in England, and best known throughout the world. It is said to have travelled from Persia to England less than two hundred years ago and to have spread from thence to other countries visited by English ships.
Though the assumptions surrounding this species' origins were not yet the same as modern ones, by the 20th century, it was believed among naturalists that the brown rat did not originate in Norway, rather the species came from central Asia and (likely) China.
Description
The fur is usually brown or dark grey, while the underparts are lighter grey or brown. The brown rat is a rather large murid and can weigh twice as much as a black rat (Rattus rattus) and many times more than a house mouse (Mus musculus). The head and body length ranges from while the tail ranges in length from , therefore being shorter than the head and body. Adult weight ranges from . Large individuals can reach but are not expected outside of domestic specimens. Stories of rats attaining sizes as big as cats are exaggerations, or misidentifications of larger rodents, such as the coypu and muskrat. It is common for breeding wild brown rats to weigh (sometimes considerably) less than . The heaviest live Rattus norvegicus on record is and they can reach a maximum length of .
Brown rats have acute hearing, are sensitive to ultrasound, and possess a very highly developed olfactory sense. Their average heart rate is 300 to 400 beats per minute, with a respiratory rate of around 100 per minute. The vision of a pigmented rat is poor, around 20/600, while a non-pigmented (albino) with no melanin in its eyes has both around 20/1200 vision and a terrible scattering of light within its vision. Brown rats are dichromats which perceive colors rather like a human with red-green colorblindness, and their colour saturation may be quite faint. Their blue perception, however, also has UV receptors, allowing them to see ultraviolet lights that humans and some other species cannot.
Biology and behavior
The brown rat is nocturnal and is a good swimmer, both on the surface and underwater, and has been observed climbing slim round metal poles several feet in order to reach garden bird feeders. Brown rats dig well, and often excavate extensive burrow systems. A 2007 study found brown rats to possess metacognition, a mental ability previously only found in humans and some other primates, but further analysis suggested they may have been following simple operant conditioning principles.
Communication
Brown rats are capable of producing ultrasonic vocalizations. As pups, young rats use different types of ultrasonic cries to elicit and direct maternal search behavior, as well as to regulate their mother's movements in the nest. Although pups produce ultrasounds around any other rats at the age of 7 days, by 14 days old they significantly reduce ultrasound production around male rats as a defensive response. Adult rats will emit ultrasonic vocalizations in response to predators or perceived danger; the frequency and duration of such cries depends on the sex and reproductive status of the rat. The female rat also emit ultrasonic vocalizations during mating.
Rats may also emit short, high frequency, ultrasonic, socially induced vocalization during rough and tumble play, before receiving morphine, or mating, and when tickled. The vocalization, described as a distinct "chirping", has been likened to laughter, and is interpreted as an expectation of something rewarding. Like most rat vocalizations, the chirping is too high in pitch for humans to hear without special equipment. Bat detectors are often used by pet owners for this purpose.
In research studies, the chirping is associated with positive emotional feelings, and social bonding occurs with the tickler, resulting in the rats becoming conditioned to seek the tickling. However, as the rats age, the tendency to chirp appears to decline.
Brown rats also produce communicative noises capable of being heard by humans. The most commonly heard in domestic rats is bruxing, or teeth-grinding, which is most usually triggered by happiness, but can also be 'self-comforting' in stressful situations, such as a visit to the vet. The noise is best described as either a quick clicking or 'burring' sound, varying from animal to animal. Vigorous bruxing can be accompanied by boggling, where the eyes of the rat rapidly bulge and retract due to movement of the lower jaw muscles behind the eye socket.
In addition, they commonly squeak along a range of tones from high, abrupt pain squeaks to soft, persistent 'singing' sounds during confrontations.
Diet
The brown rat is a true omnivore and consumes almost anything, but cereals form a substantial part of its diet. The most-liked foods of brown rats include scrambled eggs, raw carrots, and cooked corn kernels. The least-liked foods are raw beets, peaches and raw celery.
Foraging behavior is often population-specific, and varies by environment and food source. Brown rats living near a hatchery in West Virginia catch fingerling fish.
Some colonies along the banks of the Po River in Italy dive for mollusks, a practice demonstrating social learning among members of this species. Rats on the island of Norderoog in the North Sea stalk and kill sparrows and ducks.
Also preyed upon by brown rats are chicks, mice and small lizards. Examination of a wild brown rat stomachs in Germany revealed 4,000 food items, most of which were plants, although studies have shown that brown rats prefer meat when given the option. In metropolitan areas, they survive mainly on discarded human food and anything else that can be eaten without negative consequences.
Reproduction and life cycle
The brown rat can breed throughout the year if conditions are suitable, with a female producing up to five litters a year. The gestation period is only 21 days, and litters can number up to 14, although seven is common. They reach sexual maturity in about five weeks. Under ideal conditions (for the rat), this means that the population of females could increase by a factor of three and a half (half a litter of 7) in 8 weeks (5 weeks for sexual maturity and 3 weeks of gestation), corresponding to a population growing by a factor of 10 in just 15 weeks. As a result, the population can grow from 2 to 15,000 in a year. The maximum life span is three years, although most barely manage one. A yearly mortality rate of 95% is estimated, with predators and interspecies conflict as major causes.
When lactating, female rats display a 24-hour rhythm of maternal behavior, and will usually spend more time attending to smaller litters than large ones.
Brown rats live in large, hierarchical groups, either in burrows or subsurface places, such as sewers and cellars. When food is in short supply, the rats lower in social order are the first to die. If a large fraction of a rat population is exterminated, the remaining rats will increase their reproductive rate, and quickly restore the old population level.
The female is capable of becoming pregnant immediately after giving birth, and can nurse one litter while pregnant with another. She is able to produce and raise two healthy litters of normal size and weight without significantly changing her own food intake. However, when food is restricted, she can extend pregnancy by over two weeks, and give birth to litters of normal number and weight.
Mating behaviors
Males can ejaculate multiple times in a row, and this increases the likelihood of pregnancy as well as decreases the number of stillborns. Multiple ejaculation also means that males can mate with multiple females, and they exhibit more ejaculatory series when there are several oestrous females present. Males also copulate at shorter intervals than females. In group mating, females often switch partners.
Dominant males have higher mating success and also provide females with more ejaculate, and females are more likely to use the sperm of dominant males for fertilization.
In mating, female rats show a clear mating preference for unknown males versus males that they have already mated with (also known as the Coolidge effect), and will often resume copulatory behavior when introduced to a novel sexual partner.
Females also prefer to mate with males who have not experienced social stress during adolescence, and can determine which males were stressed even without any observed difference in sexual performance of males experiencing stress during adolescence and not.
Social behavior
Rats commonly groom each other and sleep together. Rats are said to establish an order of hierarchy, so one rat will be dominant over another one. Groups of rats tend to "play fight", which can involve any combination of jumping, chasing, tumbling, and "boxing". Play fighting involves rats going for each other's necks, while serious fighting involves strikes at the others' back ends. If living space becomes limited, rats may turn to aggressive behavior, which may result in the death of some animals, reducing the burden over the living space.
Rats, like most mammals, also form family groups of a mother and her young. This applies to both groups of males and females. However, rats are territorial animals, meaning that they usually act aggressively towards or scared of strange rats. Rats will fluff up their hair, hiss, squeal, and move their tails around when defending their territory. Rats will chase each other, groom each other, sleep in group nests, wrestle with each other, have dominance squabbles, communicate, and play in various other ways with each other. Huddling is an additional important part of rat socialization. Huddling, an extreme form of herding and like chattering or "bruxing" is often used to communicate that they are feeling threatened and not to come near. The common rat has been more successful at inhabiting and building communities on 6 continents and are the only species to have occupied more land than humans.
During the wintry months, rats will huddle into piles – usually cheek-to-cheek – to control humidity and keep the air warm as a heat-conserving function. Just like elderly rats are commonly groomed and nursed by their companions, nestling rats especially depend on heat from their mother, since they cannot regulate their own temperature. Other forms of interaction include: crawling under, which is literally the act of crawling underneath one another (this is common when the rat is feeling ill and helps them breathe); walking over to find a space next to their closest friend, also explained in the name; allo-grooming, so-called to distinguish it from self-grooming; and nosing, where a rat gently pushes with its nose at another rat near the neck.
Burrowing
Rats are known to burrow extensively, both in the wild and in captivity, if given access to a suitable substrate. Rats generally begin a new burrow adjacent to an object or structure, as this provides a sturdy "roof" for the section of the burrow nearest to the ground's surface. Burrows usually develop to eventually include multiple levels of tunnels, as well as a secondary entrance. Older male rats will generally not burrow, while young males and females will burrow vigorously.
Burrows provide rats with shelter and food storage, as well as safe, thermo-regulated nest sites. Rats use their burrows to escape from perceived threats in the surrounding environment; for example, rats will retreat to their burrows following a sudden, loud noise or while fleeing an intruder. Burrowing can therefore be described as a "pre-encounter defensive behavior", as opposed to a "post-encounter defensive behavior", such as flight, freezing, or avoidance of a threatening stimulus.
Distribution and habitat
Possibly originating from the plains of northern China and Mongolia, the brown rat spread to other parts of the world sometime in the Middle Ages. The question of when brown rats became commensal with humans remains unsettled, but as a species, they have spread and established themselves along routes of human migration and now live almost everywhere humans are.
The brown rat may have been present in Europe as early as 1553, a conclusion drawn from an illustration and description by Swiss naturalist Conrad Gesner in his book Historiae animalium, published 1551–1558. Though Gesner's description could apply to the black rat, his mention of a large percentage of albino specimens—not uncommon among wild populations of brown rats—adds credibility to this conclusion. Reliable reports dating to the 18th century document the presence of the brown rat in Ireland in 1722, England in 1730, France in 1735, Germany in 1750, and Spain in 1800, becoming widespread during the Industrial Revolution. It did not reach North America until around 1750–1755.
As it spread from Asia, the brown rat generally displaced the black rat in areas where humans lived. In addition to being larger and more aggressive, the change from wooden structures and thatched roofs to bricked and tiled buildings favored the burrowing brown rats over the arboreal black rats. In addition, brown rats eat a wider variety of foods, and are more resistant to weather extremes.
In the absence of humans, brown rats prefer damp environments, such as river banks. However, the great majority are now linked to man-made environments, such as sewage systems.
It is often said that there are as many rats in cities as people, but this varies from area to area depending on climate, living conditions, etc. Brown rats in cities tend not to wander extensively, often staying within of their nest if a suitable concentrated food supply is available, but they will range more widely where food availability is lower. It is difficult to determine the extent of their home range because they do not utilize a whole area but rather use regular runways to get from one location to another. There is great debate over the size of the population of rats in New York City, with estimates from almost 100 million rats to as few as 250,000. Experts suggest that New York is a particularly attractive place for rats because of its aging infrastructure and high poverty rates. In 2023, the city appointed Kathleen Corradi as the first Rat Czar, a position created to address the city's rat population. The position focuses on instituting policies measures to curb the population such as garbage regulation and additional rat trapping. In addition to sewers, rats are very comfortable living in alleyways and residential buildings, as there is usually a large and continuous food source in those areas.
In the United Kingdom, some figures show that the rat population has been rising, with estimations that 81 million rats reside in the UK Those figures would mean that there are 1.3 rats per person in the country. High rat populations in the UK are often attributed to the mild climate, which allow them higher survival rates during the winter. With the increase in global temperature and glacier retreat, it is estimated that brown rat populations will see an increase.
In tropical and desert regions, brown rat occurrence tends to be limited to human-modified habitats. Contiguous rat-free areas in the world include the continent of Antarctica, the Arctic, some isolated islands, the Canadian province of Alberta, and certain conservation areas in New Zealand. Most of Australia apart from the eastern and south-eastern coastal areas does not have reports of substantial rat occurrences.
Antarctica is uninhabitable by rats. The Arctic has extremely cold winters that rats cannot survive outdoors, and the human population density is extremely low, making it difficult for rats to travel from one habitation to another, although they have arrived in many coastal areas by ship. When the occasional rat infestation is found and eliminated, the rats are unable to re-infest it from an adjacent one. Isolated islands are also able to eliminate rat populations because of low human population density and the geographic distance from other rat populations.
Rats as invasive species
Many parts of the world have been populated by rats secondarily, where rats are now important invasive species that compete with and threaten local fauna. For instance, Norway rats reached North America between 1750 and 1775 and even in the early 20th century, from 1925 to 1927, 50% of ships entering the port of New York were rat infested.
Faroe Islands
The brown rat was first observed on the Faroe Islands in 1768. It is thought that the first individuals arrived on the southernmost island, Suðuroy, via the wreck of a Norwegian ship that had stranded on the Scottish Isle of Lewis on its way from Trondheim to Dublin. The drifting wreck, carrying brown rats, drifted northwards until it reached the village of Hvalba. Dispersion afterwards appears to have been fast, including all of Suðuroy within a year. In 1769, they were observed in Tórshavn on the southern part of Streymoy, and a decade later, in the villages in the northern part of this island.
From here, they crossed the strait and occupied Eysturoy during the years 1776 to 1779. In 1779, they reached Vagar. Whether the rats dispersed from the already established population in Suðuroy, or they were brought to the Faroe Islands with other ships is unknown. The Northern islands were invaded by the brown rat more than 100 years later, after Norwegians built and operated a whaling station in the village of Hvannasund on Borðoy from 1898 to 1920. From there, the brown rat spread to the neighbouring islands of Viðoy and Kunoy. A recent genomic analysis reveals three independent introductions of the invasive brown rat to the Faroe Islands.
Today the brown rat is found on seven of the eighteen Faroese islands, and is common in and around human habitations as well as in the wild. Although the brown rat is now common on all of the largest Faroese islands, only sparse information on the population is available in the literature. An investigation for infection with the spirochaete Leptospira interrogans did not find any infected animals, suggesting that Leptospira prevalence rates on the Faroe Islands may be among the lowest recorded worldwide.
Alaska
Hawadax Island (formerly known as Rat Island) in Alaska is thought to have been the first island in the Aleutians to be invaded by Norway rats (the Brown rat) when a Japanese ship went aground in the 1780s. They had a devastating effect on the native bird life. An eradication program was started in 2007 and the island was declared rat-free in June 2009.
Alberta
Alberta is the largest rat-free populated area in the world. Rat invasions of Alberta were stopped and rats were eliminated by very aggressive government rat control measures, starting during the 1950s.
The only Rattus species that is capable of surviving the climate of Alberta is the brown rat, which can only survive in the prairie region of the province, and even then must overwinter in buildings. Although it is a major agricultural area, Alberta is far from any seaport and only a portion of its eastern boundary with Saskatchewan provides a favorable entry route for rats. Brown rats cannot survive in the wild boreal forest to the north, the Rocky Mountains to the west, nor can they safely cross the semiarid High Plains of Montana to the south. The first brown rat did not reach Alberta until 1950, and in 1951, the province launched a rat-control program that included shooting, poisoning, and gassing rats, and bulldozing or burning down some rat-infested buildings. The effort was backed by legislation that required every person and every municipality to destroy and prevent the establishment of designated pests. If they failed, the provincial government could carry out the necessary measures and charge the costs to the landowner or municipality.
In the first year of the rat control program, of arsenic trioxide were spread throughout 8,000 buildings on farms along the Saskatchewan border. However, in 1953 the much safer and more effective rodenticide warfarin was introduced to replace arsenic. Warfarin is an anticoagulant that was approved as a drug for human use in 1954 and is much safer to use near humans and other large animals than arsenic. By 1960, the number of rat infestations in Alberta had dropped to below 200 per year. In 2002, the province finally recorded its first year with zero rat infestations, and from 2002 to 2007 there were only two infestations found. After an infestation of rats in the Medicine Hat landfill was found in 2012, the province's rat-free status was questioned, but provincial government rat control specialists brought in excavating machinery, dug out, shot, and poisoned 147 rats in the landfill, and no live rats were found thereafter. In 2013, the number of rat infestations in Alberta dropped to zero again. Alberta defines an infestation as two or more rats found at the same location, since a single rat cannot reproduce. About a dozen single rats enter Alberta in an average year and are killed by provincial rat control specialists before they can reproduce.
Only zoos, universities, and research institutes are allowed to keep caged rats in Alberta, and possession of unlicensed rats, including fancy rats by anyone else is punishable by a penalty of up to C$5,000 or up to 60 days in jail.
The adjacent and similarly landlocked province of Saskatchewan initiated a rat control program in 1972, and has managed to reduce the number of rats in the province substantially, although they have not been eliminated. The Saskatchewan rat control program has considerably reduced the number of rats trying to enter Alberta.
New Zealand
First arriving before 1800 (perhaps on James Cook's vessels), brown rats pose a serious threat to many of New Zealand's native wildlife. Rat eradication programmes within New Zealand have led to rat-free zones on offshore islands and even on fenced "ecological islands" on the mainland. Before an eradication effort was launched in 2001, the sub-Antarctic Campbell Island had the highest population density of brown rats in the world.
Diseases
Similar to other rodents, brown rats may carry a number of pathogens, which can result in disease, including Weil's disease, rat bite fever, cryptosporidiosis, viral hemorrhagic fever, Q fever and hantavirus pulmonary syndrome. In the United Kingdom, brown rats are an important reservoir for Coxiella burnetii, the bacterium that causes Q fever, with seroprevalence for the bacteria found to be as high as 53% in some wild populations.
This species can also serve as a reservoir for Toxoplasma gondii, the parasite that causes toxoplasmosis, though the disease usually spreads from rats to humans when domestic cats feed on infected brown rats. The parasite has a long history with the brown rat, and there are indications that the parasite has evolved to alter an infected rat's perception to cat predation, making it more susceptible to predation and increasing the likelihood of transmission.
Surveys and specimens of brown rat populations throughout the world have shown this species is often associated with outbreaks of trichinosis, but the extent to which the brown rat is responsible in transmitting Trichinella larvae to humans and other synanthropic animals is at least somewhat debatable. Trichinella pseudospiralis, a parasite previously not considered to be a potential pathogen in humans or domestic animals, has been found to be pathogenic in humans and carried by brown rats.
They can also be responsible for transmitting Angiostrongylus larvae to humans by eating raw or undercooked snails, slugs, molluscs, crustaceans, water and/or vegetables contaminated with them.
Brown rats are sometimes mistakenly thought to be a major reservoir of bubonic plague, a possible cause of the Black Death. However, the bacterium responsible, Yersinia pestis, is commonly endemic in only a few rodent species and is usually transmitted zoonotically by rat fleas—common carrier rodents today include ground squirrels and wood rats. However, brown rats may suffer from plague, as can many nonrodent species, including dogs, cats, and humans. During investigations of the plague epidemic in San Francisco in 1907, >1% of collected rats were infected with Y. pestis. The original carrier for the plague-infected fleas thought to be responsible for the Black Death was the black rat, and it has been hypothesized that the displacement of black rats by brown rats led to the decline of bubonic plague. This theory has, however, been deprecated, as the dates of these displacements do not match the increases and decreases in plague outbreaks.
During the COVID-19 pandemic, one study of New York City sewer rats showed that 17 percent of the city's brown rat population had become infected with SARS-CoV-2.
In captivity
Uses in science
Selective breeding of white-marked rats rescued from being killed in a now-outlawed sport called rat baiting has produced the pink-eyed white laboratory rat. Like mice, these rats are frequently subjects of medical, psychological and other biological experiments, and constitute an important model organism. This is because they grow quickly to sexual maturity and are easy to keep and to breed in captivity. When modern biologists refer to "rats", they almost always mean Rattus norvegicus.
As pets
The brown rat is kept as a pet in many parts of the world. Australia, the United Kingdom, and the United States are just a few of the countries that have formed fancy rat associations similar in nature to the American Kennel Club, establishing standards, orchestrating events, and promoting responsible pet ownership.
The many different types of domesticated brown rats include variations in coat patterns, as well as the style of the coat, such as Hairless or Rex, and more recently developed variations in body size and structure, including dwarf and tailless fancy rats.
Working rats
A working rat is a rat trained for specific tasks as a working animal. In many cases, working rats are domesticated brown rats. However, other species, notably the Gambian pouched rat, have been trained to assist humans.
| Biology and health sciences | Rodents | null |
602919 | https://en.wikipedia.org/wiki/Physical%20examination | Physical examination | In a physical examination, medical examination, clinical examination, or medical checkup, a medical practitioner examines a patient for any possible medical signs or symptoms of a medical condition. It generally consists of a series of questions about the patient's medical history followed by an examination based on the reported symptoms. Together, the medical history and the physical examination help to determine a diagnosis and devise the treatment plan. These data then become part of the medical record.
Types
Routine
The routine physical, also known as general medical examination, periodic health evaluation, annual physical, comprehensive medical exam, general health check, preventive health examination, medical check-up, or simply medical, is a physical examination performed on an asymptomatic patient for medical screening purposes. These are normally performed by a pediatrician, family practice physician, a physical therapist, physician assistant, a certified nurse practitioner or other primary care provider. This routine physical exam usually includes the HEENT evaluation. Nursing professionals such as Registered Nurse, Licensed Practical Nurses can develop a baseline assessment to identify normal versus abnormal findings. These are reported to the primary care provider. If necessary, the patient may be sent to a medical specialist for further, more detailed examinations.
The term is generally not meant to include visits for the purpose of newborn checks, Pap smears for cervical cancer, or regular visits for people with certain chronic medical disorders (for example, diabetes). The general medical examination generally involves a medical history, a (brief or complete) physical examination and sometimes laboratory tests. Some more advanced tests include ultrasound and mammography.
If done for a group of people the routine physical is a form of screening, as the aim of the examination is to detect early signs of diseases to prevent them.
Evidence
Although annual medical examinations are a routine practice in several countries, examinations performed on an asymptomatic patient are poorly supported by scientific evidence in the majority of the population. A Cochrane Collaboration meta-study found that routine annual physicals did not measurably reduce the risk of illness or death, and conversely, could lead to overdiagnosis and over-treatment; however, this article does not conclude that being in regular communication with a doctor is not important, simply that an actual physical examination may not be necessary.
Some notable general health organisations recommend against annual examinations, and propose a frequency adapted to age and previous examination results (risk factors). The specialist American Cancer Society recommends a cancer-related health check-up annually in men and women older than 40, and every three years for those older than 20.
A systematic review of studies until September 2006 concluded that the examination does result in better delivery of some other screening interventions (such as Pap smears, cholesterol screening, and faecal occult blood tests) and less patient worry. Evidence supports several of these individual screening interventions. The effects of annual check-ups on overall costs, patient disability and mortality, disease detection, and intermediate end points such a blood pressure or cholesterol, are inconclusive. A recent study found that the examination is associated with increased participation in cancer screening.
Some employers require a mandatory health checkup before hiring a candidate, even though it is now well known that some of the components of the prophylactic annual visit may actually cause harm. For example, lab tests and exams that are performed on healthy patients (as opposed to people with symptoms or known illnesses) are statistically more likely to be "false positives"—that is, when test results suggest a problem that does not exist. Disadvantages cited include the time and money that could be saved by targeted screening (health economics argument), increased anxiety over health risks (medicalisation), overdiagnosis, wrong diagnosis (for example athletic heart syndrome misdiagnosed as hypertrophic cardiomyopathy) and harm, or even death, resulting from unnecessary testing to detect or confirm, often non-existent, medical problems or while performing routine procedures as a followup after screening.
The lack of good evidence contrasts with population surveys showing that the general public is fond of these examinations, especially when they are free of charge. Despite guidelines recommending against routine annual examinations, many family physicians perform them. A fee-for-service healthcare system has been suggested to promote this practice. An alternative would be to tailor the screening interval to the age, sex, medical conditions and risk factors of each patient. This means choosing between a wide variety of tests.
Prevalence
The routine physical is commonly performed in the United States and Japan, whereas the practice varies among South East Asia and mainland European countries. In Japan it is required by law for regular working employees to have a health check once a year.
History
The roots of the periodic medical examination are not entirely clear. They have been referenced as early as 1671. They have also been advocated for since the 1920s. Some authors point to pleads from the 19th and early 20th century for the early detection of diseases like tuberculosis, and periodic school health examinations. The advent of medical insurance and related commercial influences seems to have promoted the examination, whereas this practice has been subject to controversy in the age of evidence-based medicine. Several studies have been performed before current evidence-based recommendation for screening were formulated, limiting the applicability of these studies to current-day practice.
Comprehensive
Comprehensive physical exams, also known as executive physicals, typically include laboratory tests, chest x-rays, pulmonary function testing, audiograms, full body CAT scanning, EKGs, heart stress tests, vascular age tests, urinalysis, and mammograms or prostate exams depending on gender.
Pre-employment
Pre-employment examinations are screening tests which judge the suitability of a worker for hire based on the results of their physical examination. This is also called pre-employment medical clearance. Some employers believe that by only hiring workers whose physical examination results pass certain exclusionary criteria, their employees collectively will have fewer absences due to sickness, fewer workplace injuries, and less occupational disease. A small amount of low-quality evidence in medical research supports this idea. Furthermore, the cost of staff health insurance will be lower. However, certain exams or tests that are requested by employers, such as a baseline low back x-ray, should not be performed, according to the American College of Occupational and Environmental Medicine. Reasons for this include the legality and medical necessity of the test as well as the inability of such testing to predict future problems, the radiation exposure to the worker, and the cost of the exam.
Insurance
A physical examination may be provided under health insurance cover, required of new insurance customers. This is a part of insurance medicine. In the United States, physicals are also marketed to patients as a one-stop health review, avoiding the inconvenience of attending multiple appointments with different healthcare providers.
Uses
Diagnosis
Physical examinations are performed in most healthcare encounters. For example, a physical examination is performed when a patient visits complaining of flu-like symptoms. These diagnostic examinations usually focus on the patient's chief complaint.
Screening
General health checks, including physical examinations performed when the patient reported no health concerns, often include medical screening for common conditions, such as high blood pressure. A Cochrane review found that general health checks did not reduce the risk of death from cancer, heart disease, or any other cause, and could not be proved to affect the patient's likelihood of being admitted to the hospital, becoming disabled, missing work, or needing additional office visits. The study found no effect on the risk of illness, but did find evidence suggesting that patients subject to routine physicals were diagnosed with hypertension and other chronic conditions at a higher rate than those who were not. Its authors noted that studies often failed to consider or report possible harmful outcomes (such as unwarranted anxiety or unnecessary follow-up procedures), and concluded that routine health checks were "unlikely to be beneficial" in regards to lowering cardiovascular and cancer morbidity and mortality.
Doctor-patient relations
Physical examination has been described as a ritual that plays a significant role in the doctor-patient relationship that will provide benefits in other medical encounters. When a physical exam is expected by the patient but is not performed by the provider, patients may express concern for the lack of depth of investigation into their illness, the validity of treatment plans and exclusions, and the doctor-patient relationship.
Other uses
By extension, the term "health check" is also used for routing checks on the working of equipment or business operations or solvency.
Format and interpretation
A physical examination may include checking vital signs, including temperature examination, blood pressure, pulse, and respiratory rate. The healthcare provider uses the senses of sight, hearing, touch, and sometimes smell (e.g., in infection, uremia, diabetic ketoacidosis). Taste has been made redundant by the availability of modern lab tests. Four actions are taught as the basis of physical examination: inspection, palpation (feel), percussion (tap to determine resonance characteristics), and auscultation (listen).
Scope
Although providers have varying approaches as to the sequence of body parts, a systematic examination generally starts at the head and finishes at the extremities and includes evaluation of general patient appearance and specific organ systems. After the main organ systems have been investigated by inspection, palpation, percussion, and auscultation, specific tests may follow (such as a neurological investigation, orthopedic examination) or specific tests when a particular disease is suspected (e.g. eliciting Trousseau's sign in hypocalcemia).
While the format of examination as listed below is largely as taught and expected of students, a specialist will focus on their particular field and the nature of the problem described by the patient. Hence a cardiologist will not in routine practice undertake neurological parts of the examination other than noting that the patient is able to use all four limbs on entering the consultation room and during the consultation become aware of their hearing, eyesight, and speech. Likewise an orthopaedic surgeon will examine the affected joint, but may only briefly check the heart sounds and chest to ensure that there is not likely to be any contraindication to surgery raised by the anaesthetist. A primary care physician will also generally examine the male genitals but may leave the examination of the female genitalia to a gynecologist.
With the clues obtained during the history and physical examination the healthcare provider can now formulate a differential diagnosis, a list of potential causes of the symptoms. Specific diagnostic tests (or occasionally empirical therapy) generally confirm the cause, or shed light on other, previously overlooked, causes. The physical exam is then recorded in the medical record in a standard layout which facilitates billing and other providers later reading the notes.
While elective physical exams have become more elaborate, in routine use physical exams have become less complete. This has led to editorials in medical journals about the importance of an adequate physical examination. Physicians at Stanford University medical school have introduced a set of 25 key physical examination skills that were felt to be useful.
Recording
Depending upon the chief complaint, additional sections may be included. For example, hearing may be evaluated with a specific Weber test and Rinne test, or it may be more briefly addressed in a cranial nerve exam. To give another example, a neurological related complaint might be evaluated with a specific test, such as the Romberg maneuver.
History
The Old Testament makes provision for persons in the Israelite community with leprosy to be examined by a priest: if the presenting sore was white and appeared to go beyond the depth of the skin, it was to be treated as a ritually defiling condition. A further examination was to take place seven days later.
The medical history and physical examination were supremely important to diagnosis before advanced health technology was developed, and even today, despite advances in medical imaging and molecular medical tests, the history and physical remain indispensable steps in evaluating any patient. Before the 19th century, the history and physical examination were nearly the only diagnostic tools the physician had, which explains why tactile skill and ingenious appreciation in the exam were so highly valued in the definition of what made for a good physician. Even as late as 1890, the world had no radiography or fluoroscopy, only early and limited forms of electrophysiologic testing, and no molecular biology as we know it today. Ever since this peak of the importance of the physical examination, reviewers have warned that clinical practice and medical education need to remain vigilant in appreciating the continuing need for physical examination and effectively teaching the skills to perform it; this call is ongoing, as the 21st-century literature shows.
Society and culture
People may request modesty in medical settings when the health care provider examines them.
In many Western societies, a physical exam is required to participate in extracurricular sporting activities. During the physical examination, the doctor will examine the genitals, including the penis and testicles. The doctor may ask the teenager to cough while examining the scrotum. Although this can be embarrassing for an adolescent male, it is necessary to help evaluate the presence of inguinal hernias or tumors.
| Biology and health sciences | Medical procedures | null |
602960 | https://en.wikipedia.org/wiki/Food%20processing | Food processing | Food processing is the transformation of agricultural products into food, or of one form of food into other forms. Food processing takes many forms, from grinding grain into raw flour, home cooking, and complex industrial methods used in the making of convenience foods. Some food processing methods play important roles in reducing food waste and improving food preservation, thus reducing the total environmental impact of agriculture and improving food security.
The Nova classification groups food according to different food processing techniques.
Primary food processing is necessary to make most foods edible while secondary food processing turns ingredients into familiar foods, such as bread. Tertiary food processing results in ultra-processed foods and has been widely criticized for promoting overnutrition and obesity, containing too much sugar and salt, too little fiber, and otherwise being unhealthful in respect to dietary needs of humans and farm animals.
Processing levels
Primary food processing
Primary food processing turns agricultural products, such as raw wheat kernels or livestock, into something that can eventually be eaten. This category includes ingredients that are produced by ancient processes such as drying, threshing, winnowing and milling grain, shelling nuts, and butchering animals for meat. It also includes deboning and cutting meat, freezing and smoking fish and meat, extracting and filtering oils, canning food, preserving food through food irradiation, and candling eggs, as well as homogenizing and pasteurizing milk.
Contamination and spoilage problems in primary food processing can lead to significant public health threats, as the resulting foods are used so widely. However, many forms of processing contribute to improved food safety and longer shelf life before the food spoils. Commercial food processing uses control systems such as hazard analysis and critical control points (HACCP) and failure mode and effects analysis (FMEA) to reduce the risk of harm.
Secondary food processing
Secondary food processing is the everyday process of creating food from ingredients that are ready to use. Baking bread, regardless of whether it is made at home, in a small bakery, or in a large factory, is an example of secondary food processing. Fermenting fish and making wine, beer, and other alcoholic products are traditional forms of secondary food processing. Sausages are a common form of secondary processed meat, formed by comminution (grinding) of meat that has already undergone primary processing. Most of the secondary food processing methods known to humankind are commonly described as cooking methods.
Tertiary food processing
Tertiary food processing is the commercial production of what is commonly called processed food. These are ready-to-eat or heat-and-serve foods, such as frozen meals and re-heated airline meals.
History
Food processing dates back to the prehistoric ages when crude processing incorporated fermenting, sun drying, preserving with salt, and various types of cooking (such as roasting, smoking, steaming, and oven baking), Such basic food processing involved chemical enzymatic changes to the basic structure of food in its natural form, as well served to build a barrier against surface microbial activity that caused rapid decay. Salt-preservation was especially common for foods that constituted warrior and sailors' diets until the introduction of canning methods. Evidence for the existence of these methods can be found in the writings of the ancient Greek, Chaldean, Egyptian and Roman civilizations as well as archaeological evidence from Europe, North and South America and Asia. These tried and tested processing techniques remained essentially the same until the advent of the Industrial Revolution. Examples of ready-meals also date back to before the preindustrial revolution, and include dishes such as Cornish pasty and Haggis. Both during ancient times and today in modern society these are considered processed foods.
Modern food processing technology developed in the 19th and 20th centuries was developed in a large part to serve military needs. In 1809, Nicolas Appert invented a hermetic bottling technique that would preserve food for French troops which ultimately contributed to the development of tinning, and subsequently canning by Peter Durand in 1810. Although initially expensive and somewhat hazardous due to the lead used in cans, canned goods would later become a staple around the world. Pasteurization, discovered by Louis Pasteur in 1864, improved the quality and safety of preserved foods and introduced the wine, beer, and milk preservation.
In the 20th century, World War II, the space race and the rising consumer society in developed countries contributed to the growth of food processing with such advances as spray drying, evaporation, juice concentrates, freeze drying and the introduction of artificial sweeteners, colouring agents, and such preservatives as sodium benzoate. In the late 20th century, products such as dried instant soups, reconstituted fruits and juices, and self cooking meals such as MRE food ration were developed. By the 20th century, automatic appliances like microwave oven, blender, and rotimatic paved way for convenience cooking.
In western Europe and North America, the second half of the 20th century witnessed a rise in the pursuit of convenience. Food processing companies marketed their products especially towards middle-class working wives and mothers. Frozen foods (often credited to Clarence Birdseye) found their success in sales of juice concentrates and "TV dinners". Processors utilised the perceived value of time to appeal to the postwar population, and this same appeal contributes to the success of convenience foods today.
Benefits and drawbacks
Benefits
Benefits of food processing include toxin removal, preservation, easing marketing and distribution tasks, and increasing food consistency. In addition, it increases yearly availability of many foods, enables transportation of delicate perishable foods across long distances and makes many kinds of foods safe to eat by de-activating spoilage and pathogenic micro-organisms. Modern supermarkets would not exist without modern food processing techniques, and long voyages would not be possible.
Processed foods are usually less susceptible to early spoilage than fresh foods and are better suited for long-distance transportation from the source to the consumer. When they were first introduced, some processed foods helped to alleviate food shortages and improved the overall nutrition of populations as it made many new foods available to the masses.
Processing can also reduce the incidence of food-borne disease. Fresh materials, such as fresh produce and raw meats, are more likely to harbour pathogenic micro-organisms (e.g. Salmonella) capable of causing serious illnesses.
The extremely varied modern diet is only truly possible on a wide scale because of food processing. Transportation of more exotic foods, as well as the elimination of much hard labor gives the modern eater easy access to a wide variety of food unimaginable to their ancestors.
The act of processing can often improve the taste of food significantly.
Mass production of food is much cheaper overall than individual production of meals from raw ingredients. Therefore, a large profit potential exists for the manufacturers and suppliers of processed food products. Individuals may see a benefit in convenience, but rarely see any direct financial cost benefit in using processed food as compared to home preparation.
Processed food freed people from the large amount of time involved in preparing and cooking "natural" unprocessed foods. The increase in free time allows people much more choice in life style than previously allowed. In many families the adults are working away from home and therefore there is little time for the preparation of food based on fresh ingredients. The food industry offers products that fulfill many different needs: e.g. fully prepared ready meals that can be heated up in the microwave oven within a few minutes.
Modern food processing also improves the quality of life for people with allergies, diabetics, and other people who cannot consume some common food elements. Food processing can also add extra nutrients such as vitamins.
Drawbacks
Processing of food can decrease its nutritional density. The amount of nutrients lost depends on the food and processing method. For example, heat destroys vitamin C. Therefore, canned fruits possess less vitamin C than their fresh alternatives. The USDA conducted a study of nutrient retention in 2004, creating a table of foods, levels of preparation, and nutrition.
New research highlighting the importance to human health of a rich microbial environment in the intestine indicates that abundant food processing (not fermentation of foods) endangers that environment.
Added sodium
One of the main sources for sodium in the diet is processed foods. Sodium, mostly in the form of sodium chloride, i.e. salt, is added to prevent spoilage, add flavor and enhance the texture of these foods. Americans consume an average of 3436 milligrams of sodium per day, which is higher than the recommended limit of 2300 milligrams per day for healthy people, and more than twice the limit of 1500 milligrams per day for those at increased risk for heart disease.
Added sugars
While it is not necessary to limit the sugars found naturally in whole, unprocessed foods like fresh fruit, eating too much added sugar found in many processed foods increases the risk of heart disease, obesity, cavities and Type 2 diabetes. The American Heart Association recommends women limit added sugars to no more than , or 25 grams, and men limit added sugars to no more than , or about 38.75 grams, per day. Currently, Americans consume an average of from added sugars each day.
Nutrient losses
Processing foods often involves nutrient losses, which can make it harder to meet the body's needs if these nutrients are not added back through fortification or enrichment. For example, using high heat during processing can cause vitamin C losses. Another example is refined grains, which have less fiber, vitamins and minerals than whole grains. Eating refined grains, such as those found in many processed foods, instead of whole grains may increase the risk for high cholesterol, diabetes and obesity, according to a study published in "The American Journal of Clinical Nutrition" in December 2007.
Trans fats
Foods that have undergone processing, including some commercial baked goods, desserts, margarine, frozen pizza, microwave popcorn and coffee creamers, sometimes contain trans fats. This is the most unhealthy type of fat, and may increase risk for high cholesterol, heart disease and stroke. The 2010 Dietary Guidelines for Americans recommends keeping trans fat intake as low as possible.
Other potential disadvantages
Processed foods may actually take less energy to digest than whole foods, according to a study published in "Food & Nutrition Research" in 2010, meaning more of their food energy content is retained within the body. Processed foods also tend to be more allergenic than whole foods, according to a June 2004 "Current Opinion in Allergy and Clinical Immunology" article. Although the preservatives and other food additives used in many processed foods are generally recognized as safe, a few may cause problems for some individuals, including sulfites, artificial sweeteners, artificial colors and flavors, sodium nitrate, BHA and BHT, olestra, caffeine and monosodium glutamate — a flavor enhancer.
Performance parameters for food processing
When designing processes for the food industry the following performance parameters may be taken into account:
Hygiene, e.g. measured by number of micro-organisms per mL of finished product.
Energy efficiency measured e.g. by "ton of steam per ton of sugar produced".
Minimization of waste, measured e.g. by "percentage of peeling loss during the peeling of potatoes".
Labour used, measured e.g. by "number of working hours per ton of finished product".
Minimization of cleaning stops measured e.g. by "number of hours between cleaning stops".
Industries
Food processing industries and practices include the following:
Cannery
Fish processing
Food packaging plant
Industrial rendering
Meat packing plant
Potato processing industry
Slaughterhouse
Sugar industry
| Technology | Food, water and health | null |
603273 | https://en.wikipedia.org/wiki/Magnification | Magnification | Magnification is the process of enlarging the apparent size, not physical size, of something. This enlargement is quantified by a size ratio called optical magnification. When this number is less than one, it refers to a reduction in size, sometimes called de-magnification.
Typically, magnification is related to scaling up visuals or images to be able to see more detail, increasing resolution, using microscope, printing techniques, or digital processing. In all cases, the magnification of the image does not change the perspective of the image.
Examples of magnification
Some optical instruments provide visual aid by magnifying small or distant subjects.
A magnifying glass, which uses a positive (convex) lens to make things look bigger by allowing the user to hold them closer to their eye.
A telescope, which uses its large objective lens or primary mirror to create an image of a distant object and then allows the user to examine the image closely with a smaller eyepiece lens, thus making the object look larger.
A microscope, which makes a small object appear as a much larger image at a comfortable distance for viewing. A microscope is similar in layout to a telescope except that the object being viewed is close to the objective, which is usually much smaller than the eyepiece.
A slide projector, which projects a large image of a small slide on a screen. A photographic enlarger is similar.
A zoom lens, a system of camera lens elements for which the focal length and angle of view can be varied.
Size ratio (optical magnification)
Optical magnification is the ratio between the apparent size of an object (or its size in an image) and its true size, and thus it is a dimensionless number. Optical magnification is sometimes referred to as "power" (for example "10× power"), although this can lead to confusion with optical power.
Linear or transverse magnification
For real images, such as images projected on a screen, size means a linear dimension (measured, for example, in millimeters or inches).
Angular magnification
For optical instruments with an eyepiece, the linear dimension of the image seen in the eyepiece (virtual image at infinite distance) cannot be given, thus size means the angle subtended by the object at the focal point (angular size). Strictly speaking, one should take the tangent of that angle (in practice, this makes a difference only if the angle is larger than a few degrees). Thus, angular magnification is given by:
where is the angle subtended by the object at the front focal point of the objective and is the angle subtended by the image at the rear focal point of the eyepiece.
For example, the mean angular size of the Moon's disk as viewed from Earth's surface is about 0.52°. Thus, through binoculars with 10× magnification, the Moon appears to subtend an angle of about 5.2°.
By convention, for magnifying glasses and optical microscopes, where the size of the object is a linear dimension and the apparent size is an angle, the magnification is the ratio between the apparent (angular) size as seen in the eyepiece and the angular size of the object when placed at the conventional closest distance of distinct vision: from the eye.
By instrument
Single lens
The linear magnification of a thin lens is
where is the focal length, is the distance from the lens to the object, and as the distance of the object with respect to the front focal point. A sign convention is used such that and (the image distance from the lens) are positive for real object and image, respectively, and negative for virtual object and images, respectively. of a converging lens is positive while for a diverging lens it is negative.
For real images, is negative and the image is inverted. For virtual images, is positive and the image is upright.
With being the distance from the lens to the image, the height of the image and the height of the object, the magnification can also be written as:
Note again that a negative magnification implies an inverted image.
The image magnification along the optical axis direction , called longitudinal magnification, can also be defined. The Newtonian lens equation is stated as , where and as on-axis distances of an object and the image with respect to respective focal points, respectively. is defined as
and by using the Newtonian lens equation,
The longitudinal magnification is always negative, means that, the object and the image move toward the same direction along the optical axis. The longitudinal magnification varies much faster than the transverse magnification, so the 3-dimensional image is distorted.
Photography
The image recorded by a photographic film or image sensor is always a real image and is usually inverted. When measuring the height of an inverted image using the cartesian sign convention (where the x-axis is the optical axis) the value for will be negative, and as a result will also be negative. However, the traditional sign convention used in photography is "real is positive, virtual is negative". Therefore, in photography: Object height and distance are always and positive. When the focal length is positive the image's height, distance and magnification are and positive. Only if the focal length is negative, the image's height, distance and magnification are and negative. Therefore, the formulae are traditionally presented as
Magnifying glass
The maximum angular magnification (compared to the naked eye) of a magnifying glass depends on how the glass and the object are held, relative to the eye. If the lens is held at a distance from the object such that its front focal point is on the object being viewed, the relaxed eye (focused to infinity) can view the image with angular magnification
Here, is the focal length of the lens in centimeters. The constant 25 cm is an estimate of the "near point" distance of the eye—the closest distance at which the healthy naked eye can focus. In this case the angular magnification is independent from the distance kept between the eye and the magnifying glass.
If instead the lens is held very close to the eye and the object is placed closer to the lens than its focal point so that the observer focuses on the near point, a larger angular magnification can be obtained, approaching
A different interpretation of the working of the latter case is that the magnifying glass changes the diopter of the eye (making it myopic) so that the object can be placed closer to the eye resulting in a larger angular magnification.
Microscope
The angular magnification of a microscope is given by
where is the magnification of the objective and the magnification of the eyepiece. The magnification of the objective depends on its focal length and on the distance between objective back focal plane and the focal plane of the eyepiece (called the tube length):
The magnification of the eyepiece depends upon its focal length and is calculated by the same equation as that of a magnifying glass:
Note that both astronomical telescopes as well as simple microscopes produce an inverted image, thus the equation for the magnification of a telescope or microscope is often given with a minus sign.
Telescope
The angular magnification of an optical telescope is given by
in which is the focal length of the objective lens in a refractor or of the primary mirror in a reflector, and is the focal length of the eyepiece.
Measurement of telescope magnification
Measuring the actual angular magnification of a telescope is difficult, but it is possible to use the reciprocal relationship between the linear magnification and the angular magnification, since the linear magnification is constant for all objects.
The telescope is focused correctly for viewing objects at the distance for which the angular magnification is to be determined and then the object glass is used as an object the image of which is known as the exit pupil. The diameter of this may be measured using an instrument known as a Ramsden dynameter which consists of a Ramsden eyepiece with micrometer hairs in the back focal plane. This is mounted in front of the telescope eyepiece and used to evaluate the diameter of the exit pupil. This will be much smaller than the object glass diameter, which gives the linear magnification (actually a reduction), the angular magnification can be determined from
Maximum usable magnification
With any telescope, microscope or lens,
a maximum magnification exists beyond which the image looks bigger but shows no more detail. It occurs when the finest detail the instrument can resolve is magnified to match the finest detail the eye can see. Magnification beyond this maximum is sometimes called "empty magnification".
For a good quality telescope operating in good atmospheric conditions, the maximum usable magnification is limited by diffraction. In practice it is considered to be 2× the aperture in millimetres or 50× the aperture in inches; so, a diameter telescope has a maximum usable magnification of 120×.
With an optical microscope having a high numerical aperture and using oil immersion, the best possible resolution is corresponding to a magnification of around 1200×. Without oil immersion, the maximum usable magnification is around 800×. For details, see limitations of optical microscopes.
Small, cheap telescopes and microscopes are sometimes supplied with the eyepieces that give magnification far higher than is usable.
The maximum relative to the minimum magnification of an optical system is known as zoom ratio.
"Magnification" of displayed images
Magnification figures on pictures displayed in print or online can be misleading. Editors of journals and magazines routinely resize images to fit the page, making any magnification number provided in the figure legend incorrect. Images displayed on a computer screen change size based on the size of the screen. A scale bar (or micron bar) is a bar of stated length superimposed on a picture. When the picture is resized the bar will be resized in proportion. If a picture has a scale bar, the actual magnification can easily be calculated. Where the scale (magnification) of an image is important or relevant, including a scale bar is preferable to stating magnification.
| Physical sciences | Optics | Physics |
604156 | https://en.wikipedia.org/wiki/Nicobar%20pigeon | Nicobar pigeon | The Nicobar pigeon or Nicobar dove (Caloenas nicobarica, Car: ) is a bird found on small islands and in coastal regions from the Andaman and Nicobar Islands, India, east through the Indonesian Archipelago, to the Solomons and Palau. It is the only living member of the genus Caloenas alongside the extinct spotted green pigeon and Kanaka pigeon, and is the closest living relative of the extinct dodo and Rodrigues solitaire.
Taxonomy
In 1738, the English naturalist Eleazar Albin included a description and two illustrations of the Nicobar pigeon in his A Natural History of Birds. When in 1758 the Swedish naturalist Carl Linnaeus updated his Systema Naturae for the tenth edition, he placed the Nicobar pigeon with all the other pigeons in the genus Columba. Linnaeus included a brief description, coined the binomial name Columba nicobarica and cited Albin's work. The species is now placed in the genus Caloenas erected by English zoologist George Robert Gray in 1840 with the Nicobar pigeon as the type species.
Two subspecies are recognised:
C. n. nicobarica (Linnaeus, 1758) – Andaman and Nicobar Islands, Malay Archipelago to New Guinea, Philippines and Solomon Islands
C. n. pelewensis Finsch, 1875 – Palau Island
Based on cladistic analysis of mtDNA cytochrome b and 12S rRNA sequences, the Nicobar pigeon is sometimes called the closest living relative of the extinct didines (Raphinae), which include the famous dodo (Raphus cucullatus). However, the study's results showed this as one weak possibility from a limited sample of taxa. In any case, nDNA β-fibrinogen intron 7 sequence data agrees with the idea of the Raphinae as a subfamily of pigeons (and not an independent family, as was previously believed due to their bizarre apomorphies) that was part of a diverse Indopacific radiation, to which the Nicobar pigeon also belongs.
The following cladogram, from Shapiro and colleagues (2002), shows the Nicobar pigeon's closest relationships within Columbidae, a clade consisting of generally ground-dwelling island endemics.
A similar cladogram was published in 2007, differing only in the inverted placement of Goura and Didunculus, as well as in the inclusion of the pheasant pigeon and the thick-billed ground pigeon at the base of the clade.
C. nicobarica is a quite singular columbiform (though less autapomorphic than the flightless Raphinae), as are for example the tooth-billed pigeon (Didunculus strigirostris) and the crowned pigeons (Goura), which are typically considered distinct subfamilies. Hence, the Nicobar pigeon may well constitute another now-monotypic subfamily. And while any of the semi-terrestrial pigeons of Southeast Asia and the Wallacea cannot be excluded as possible closest living relative of the Raphinae, the Nicobar pigeon makes a more plausible candidate than for example the group of imperial-pigeons and fruit-doves, which seems to be part of the same radiation.
Whether it is possible to clarify such deep-time phylogenies without a comprehensive study of all major lineages of living Columbidae remains to be seen. The primitive molecular clock used to infer the date the ancestors of the Nicobar pigeon and the didines diverged has since turned out to be both unreliable and miscalibrated. But what little evidence is available still suggests that the Nicobar pigeon is distinct from all other living lifeforms since the Paleogene – most likely some time between 56-34 million years ago during the Eocene, which makes up the bulk of the Paleogene period.
From subfossil bones found on New Caledonia and Tonga, an extinct species of Caloenas, the Kanaka pigeon (C. canacorum) was described. It was about one-quarter larger than the Nicobar pigeon. Considering that it must have been a good source of food, it was most likely hunted to extinction by the first human settlers of its home islands. It probably was extinct by 500 BC. The Spotted green pigeon (C. maculata) is a more recently extinct species from an unknown Pacific locality; it probably disappeared in the 19th century and most likely succumbed to introduced European rats. It is placed in Caloenas as the least awkward possibility; its true affinities are presently indeterminate and it is perhaps more likely to represent a distinct genus of the Indopacific radiation of Columbidae.
Description
It is a large pigeon, measuring in length. The head is grey, like the upper neck plumage, which turns into green and copper hackles. The tail is very short and pure white. The rest of its plumage is metallic green. The cere of the dark bill forms a small blackish knob; the strong legs and feet are dull red. The irises are dark.
Females are slightly smaller than males; they have a smaller bill knob, shorter hackles and browner underparts. Immature birds have a black tail and lack almost all iridescence. There is hardly any variation across the birds' wide range. Even the Palau subspecies C. n. pelewensis has merely shorter neck hackles, but is otherwise almost identical.
It is a very vocal species, giving a low-pitched repetitive call.
Distribution and habitat
On the Nicobar Islands (which are referred to in its common and scientific names), the most significant colony of Nicobar pigeons in modern times was found on Batti Malv, a remote wildlife sanctuary between Car Nicobar and Teressa. The 2004 Indian Ocean tsunami caused massive damage on the Nicobar Islands, and to what extent Batti Malv was affected is still not clear. But while everything on some islets in the Great Nicobar Biosphere Reserve was destroyed, Batti Malv lighthouse – a skeletal tower a dozen metres high, standing a few metres ASL at the highest point of the low-lying island – was little-damaged and put back in operation by the survey ship INS Sandhayak less than one month after the disaster. An April 2007 survey by the Indian Coast Guard vessel ICGS Vikram found the lighthouse tower "totally covered" in vines, indicating rampant regeneration of vegetation – but perhaps also that damage to the island's forest was severe, as a cover of creeping plants is typical of early succession stages, while a photo of the lighthouse taken before the tsunami shows rather mature forest.
Found in Australia
A Nicobar pigeon was found by the Bardi Jawi Indigenous rangers on the Dampier Peninsula in the western Kimberley region of Australia in May 2017. As part of biosecurity measures, it was reported to quarantine services and was removed by Australian Department of Agriculture officials. In 2023, another individual was found on Green Island, off the coast of Cairns. Green Island Resort's environmental manager contacted authorities, but there are currently no plans to relocate the pigeon.
Behaviour and ecology
The Nicobar pigeon's breeding range encompasses the Andaman and Nicobar Islands of India, the Mergui Archipelago of Myanmar, offshore islands of south-western Thailand, Peninsular Malaysia, southern Cambodia and Vietnam, and many of the small islands between Sumatra, the Philippines and the Solomon Islands. On Palau, the only distinct subspecies C. n. pelewensis is found.
The Nicobar pigeon roams in flocks from island to island, usually sleeping on offshore islets where no predators occur and spends the day in areas with better food availability, not shying away from areas inhabited by humans. Its food consists of seeds, fruit and buds, and it is attracted to areas where grain is available. A gizzard stone helps to grind up hard food items. Its flight is quick, with regular beats and an occasional sharp flick of the wings, as is characteristic of pigeons in general. Unlike other pigeons, groups tend to fly in columns or single file, not in a loose flock. The white tail is prominent in flight when seen from behind and may serve as a sort of "taillight", keeping flocks together when crossing the sea at dawn or dusk. The young birds' lack of a white tail is a signal of their immaturity clearly visible to conspecifics – to an adult Nicobar pigeon, it is obvious at a glance which flockmembers are neither potential mates, nor potential competitors for mates, nor old enough to safely guide a flock from one island to another.
This species nests in dense forest on offshore islets, often in large colonies. It builds a loose stick nest in a tree. It lays one elliptical faintly blue-tinged white egg.
In 2017, several individual Caloenas nicobarica were sighted in the Kimberley region of Western Australia with a juvenile captured at Ardyaloon (One Arm Point), near Broome - the first time the bird has been sighted on the Australian mainland.
Conservation
Nicobar pigeons are hunted in considerable numbers for food, and also for their gizzard stone which is used in jewellery. The species is also trapped for the local pet market, but as it is on CITES Appendix I, commercial international trade is prohibited. Internationally, captive breeding is able to supply the birds demanded by zoos, where this attractive and unusual bird is often seen. Direct exploitation of the species, even including the illegal trade, might be sustainable on its own; however, its available nesting habitat is decreasing. The offshore islets which it requires are often logged for plantations, destroyed by construction activity, or polluted by nearby industry or harbours. Also, increased travel introduces predators to more and more of the breeding sites, and colonies of the Nicobar pigeon may be driven to desert such locations or be destroyed outright. Though the bird is widely distributed and in some locations very common --even on small Palau it is still reasonably plentiful, with an estimated 1,000 adult birds remaining—its long-term future is increasingly jeopardized. For these reasons, the IUCN considers C. nicobarica a near threatened species.
Gallery
| Biology and health sciences | Columbimorphae | Animals |
604685 | https://en.wikipedia.org/wiki/Cutibacterium%20acnes | Cutibacterium acnes | Cutibacterium acnes (Propionibacterium acnes) is the relatively slow-growing, typically aerotolerant anaerobic, gram-positive bacterium (rod) linked to the skin condition of acne; it can also cause chronic blepharitis and endophthalmitis, the latter particularly following intraocular surgery. Its genome has been sequenced and a study has shown several genes can generate enzymes for degrading skin and proteins that may be immunogenic (activating the immune system).
The species is largely commensal and part of the skin flora present on most healthy adult humans' skin. It is usually just barely detectable on the skin of healthy preadolescents. It lives, among other things, primarily on fatty acids in sebum secreted by sebaceous glands in the follicles. It may also be found throughout the gastrointestinal tract.
Originally identified as Bacillus acnes, it was later named Propionibacterium acnes for its ability to generate propionic acid. In 2016, P. acnes was taxonomically reclassified as a result of biochemical and genomic studies. In terms of both phylogenetic tree structure and DNA G + C content, the cutaneous species was distinguishable from other species that had been previously categorized as P. acnes. As part of restructuring, the novel genus Cutibacterium was created for the cutaneous species, including those formerly identified as Propionibacterium acnes, Propionibacterium avidum, and Propionibacterium granulosum. Characterization of phylotypes of C. acnes is an active field of research.
Role in diseases
Acne vulgaris
Cutibacterium acnes bacteria predominantly live deep within follicles and pores, although they are also found on the surface of healthy skin. In these follicles, C. acnes bacteria use sebum, cellular debris and metabolic byproducts from the surrounding skin tissue as their primary sources of energy and nutrients. Elevated production of sebum by hyperactive sebaceous glands (sebaceous hyperplasia) or blockage of the follicle can cause C. acnes bacteria to grow and multiply.
Cutibacterium acnes bacteria secrete many proteins, including several digestive enzymes. These enzymes are involved in the digestion of sebum and the acquisition of other nutrients. They can also destabilize the layers of cells that form the walls of the follicle. The cellular damage, metabolic byproducts and bacterial debris produced by the rapid growth of C. acnes in follicles can trigger inflammation. This inflammation can lead to the symptoms associated with some common skin disorders, such as folliculitis and acne vulgaris. Acne vulgaris is the disease most commonly associated with C. acnes infection. Cutibacterium acnes is one of the most common and universal skin diseases, affecting more than 45 million individuals in the United States. 20% of all dermatologist visits are related to treating acne-related issues. This issue often develops during hormonal periods; however, it is also apparent through early adulthood. There are no quantitative differences between C. acnes of the skin of patients with acne, but the C. acnes phylogenetic groups display distinct genetic and phenotypic characteristics. C. acnes biofilms are also found much more frequently in acne and can induce distinct immune responses to combat against acne.
Acne vulgaris is a chronic inflammatory disease of the pilosebaceous unit, which includes the hair follicle, hair shaft, and sebaceous gland and about 650 million people are affected globally by this disease. C. acnes starts to colonize on the skin around 1 to 3 years prior to puberty and grows exponentially during this time. This is why so many teens and young adults struggle with acne. Prescriptions to treat acne are often antibiotics. However, with the rise of antibiotic resistance, antibiotics are now often combined with broad-spectrum antibacterial agents such as benzoyl peroxide, and other medications like isotretinoin (commonly known by the brand name Accutane) are being used on patients with severe or resistant acne.
Staphylococcus epidermidis
The damage caused by C. acnes and the associated inflammation make the affected tissue more susceptible to colonization by opportunistic bacteria, such as Staphylococcus aureus. Preliminary research shows healthy pores are only colonized by C. acnes, while unhealthy ones universally include the nonpore-resident Staphylococcus epidermidis, amongst other bacterial contaminants. Whether this is a root causality, just opportunistic and a side effect, or a more complex pathological duality between C. acnes and this particular Staphylococcus species is not known. Current research has pointed to the idea that C. acnes and S. epidermidis have a symbiotic relationship. Both bacteria exist on the normal flora of the skin and a disrupt in balance of these bacteria on the skin can result in acne or other bacterial infection.
In addition to contributing to skin inflammation and acne lesions, an imbalance in these bacteria may also impair the skin's ability to heal and regenerate, leading to prolonged and more severe acne outbreaks. This disruption can also affect the skin's overall microbiome diversity, potentially increasing susceptibility to other skin conditions such as eczema or rosacea. Investigating the dynamics of this relationship may offer insights into novel therapeutic approaches for managing various skin disorders.
Ophthalmic complications
Cutibacterium acnes is a common cause of chronic endophthalmitis following cataract surgery. The pathogen may also cause corneal ulcers.
Disk herniation
Cutibacterium acnes has been found in herniated discs. The propionic acid which it secretes creates micro-fractures of the surrounding bone. These micro-fractures are sensitive and it has been found that antibiotics have been helpful in resolving this type of low back pain.
Sarcoidosis
Cutibacterium acnes can be found in bronchoalveolar lavage of approximately 70% of patients with sarcoidosis and is associated with disease activity, but it can also be found in 23% of controls. The subspecies of C. acnes that cause these infections of otherwise sterile tissues (prior to medical procedures), however, are the same subspecies found on the skin of individuals who do not have acne-prone skin, so are likely local contaminants. Moderate to severe acne vulgaris appears to be more often associated with virulent strains.
Opportunistic diseases
Cutibacterium acnes is often considered an opportunistic pathogen, causing a range of postoperative and device-related infections, notably e.g., surgical infections, post-neurosurgical infections, infected joint prostheses (especially shoulder), neurosurgical shunt infections and endocarditis in patients with prosthetic heart valves (predominantly men). C. acnes may play a role in other conditions, including SAPHO (synovitis, acne, pustulosis, hyperostosis, osteitis) syndrome, sarcoidosis and sciatica. It is also suspected a main bacterial source of neuroinflammation in Alzheimer's disease brains. It is a common contaminant in blood and cerebrospinal fluid cultures.
Antimicrobial susceptibility
Cutibacterium acnes bacteria are susceptible to a wide range of antimicrobial molecules, from both pharmaceutical and natural sources. The antibiotics most frequently used to treat acne vulgaris are erythromycin, clindamycin, doxycycline, and minocycline. Several other families of antibiotics are also active against C. acnes bacteria, including quinolones, cephalosporins, pleuromutilins, penicillins, and sulfonamides.
Antibiotic-resistance
The emergence of antibiotic-resistant C. acnes bacteria represents a growing problem worldwide. The problem is especially pronounced in North America and Europe. The antibiotic families that C. acnes are most likely to acquire resistance to are the macrolides (e.g., erythromycin and azithromycin), lincosamides (e.g., clindamycin) and tetracyclines (e.g., doxycycline and minocycline).
However, C. acnes bacteria are susceptible to many types of antimicrobial chemicals found in over-the-counter antibacterial products, including benzoyl peroxide, triclosan, chloroxylenol, and chlorhexidine gluconate.
C. acnes resistance to antibiotics has increased to 64% in 2000, up from 20% in 1979. Treatments such as oral macrolides are often avoided because the bacteria has become resistant in most cases. This creates a public health issue, forcing healthcare providers to seek out other forms of treatment.
Treatments
Several naturally occurring molecules and compounds are toxic to C. acnes bacteria. Some essential oils such as rosemary, tea tree oil, clove oil, and citrus oils contain antibacterial chemicals. Natural honey has also been shown to have some antibacterial properties that may be active against C. acnes.
The elements silver, sulfur, and copper have also been demonstrated to be toxic towards many bacteria, including C. acnes.
Salicylic acid is a naturally occurring substance derived from plants (white willow bark and wintergreen leaves) used to promote exfoliation of the skin in order to treat acne. Additionally, research investigates the mechanism by which salicylic acid (SA) treats acne vulgaris. A study finds that SA suppresses the AMPK/SREBP1 (AMP-activated protein kinase)(AMPK/SREBP1 pathway is a signaling pathway involved in the regulation of lipid metabolism in sebocytes, which are the cells responsible for producing sebum in the skin) pathway in sebocytes, leading to a decrease in lipid synthesis and sebum production. SA also reduces the inflammatory response of sebocytes and decreases the proliferation of C. acnes. These results suggest that SA has a multifaceted approach in treating acne vulgaris by targeting several key factors that contribute to its development. The minimum inhibitory concentration for SA is 4000–8000 μg/mL.
Photosensitivity
Cutibacterium acnes glows orange when exposed to blacklight, possibly due to the presence of endogenous porphyrins. It is also killed by ultraviolet light. C. acnes is especially sensitive to light in the 405–420 nanometer (near the ultraviolet) range due to an endogenic porphyrin–coporphyrin III. A total radiant exposure of 320 J/cm2 inactivates this species in vitro. Its photosensitivity can be enhanced by pretreatment with aminolevulinic acid, which boosts production of this chemical, although this causes significant side effects in humans, and in practice was not significantly better than the light treatment alone.
Other habitats
Cutibacterium acnes has been found to be an endophyte of plants. Notably, grapevine appears to host an endophytic population of C. acnes that is closely related to the human-associated strains. The two lines diverged roughly 7,000 years ago, at about the same time when grapevine agriculture may have been established. This C. acnes subtype was dubbed Zappae in honour of the eccentric composer Frank Zappa, to highlight its unexpected and unconventional habitat.
| Biology and health sciences | Gram-positive bacteria | Plants |
604798 | https://en.wikipedia.org/wiki/Joule%20heating | Joule heating | Joule heating (also known as resistive heating, resistance heating, or Ohmic heating) is the process by which the passage of an electric current through a conductor produces heat.
Joule's first law (also just Joule's law), also known in countries of the former USSR as the Joule–Lenz law, states that the power of heating generated by an electrical conductor equals the product of its resistance and the square of the current. Joule heating affects the whole electric conductor, unlike the Peltier effect which transfers heat from one electrical junction to another.
Joule-heating or resistive-heating is used in many devices and industrial processes. The part that converts electricity into heat is called a heating element.
Among the applications are:
Buildings are often heated with electric heaters where grid power is available.
Electric stoves and ovens use Joule heating to cook food.
Soldering irons generate heat to melt conductive solder and make electrical connections.
Cartridge heaters are used in various manufacturing processes.
Electric fuses are used as a safety device, breaking a circuit by melting if enough current flows to heat them to the melting point.
Electronic cigarettes vaporize liquid by Joule heating.
Some food processing equipment may make use of Joule heating: running a current through food material (which behave as an electrical resistor) causes heat release inside the food. The alternating electrical current coupled with the resistance of the food causes the generation of heat. A higher resistance increases the heat generated. Ohmic heating allows for fast and uniform heating of food products, which maintains quality. Products with particulates heat up faster (compared to conventional heat processing) due to higher resistance.
History
James Prescott Joule first published in December 1840, an abstract in the Proceedings of the Royal Society, suggesting that heat could be generated by an electrical current. Joule immersed a length of wire in a fixed mass of water and measured the temperature rise due to a known current flowing through the wire for a 30 minute period. By varying the current and the length of the wire he deduced that the heat produced was proportional to the square of the current multiplied by the electrical resistance of the immersed wire.
In 1841 and 1842, subsequent experiments showed that the amount of heat generated was proportional to the chemical energy used in the voltaic pile that generated the template. This led Joule to reject the caloric theory (at that time the dominant theory) in favor of the mechanical theory of heat (according to which heat is another form of energy).
Resistive heating was independently studied by Heinrich Lenz in 1842.
The SI unit of energy was subsequently named the joule and given the symbol J. The commonly known unit of power, the watt, is equivalent to one joule per second.
Microscopic description
Joule heating is caused by interactions between charge carriers (usually electrons) and the body of the conductor.
A potential difference (voltage) between two points of a conductor creates an electric field that accelerates charge carriers in the direction of the electric field, giving them kinetic energy. When the charged particles collide with the quasi-particles in the conductor (i.e. the canonically quantized, ionic lattice oscillations in the harmonic approximation of a crystal), energy is being transferred from the electrons to the lattice (by the creation of further lattice oscillations). The oscillations of the ions are the origin of the radiation ("thermal energy") that one measures in a typical experiment.
Power loss and noise
Joule heating is referred to as ohmic heating or resistive heating because of its relationship to Ohm's Law. It forms the basis for the large number of practical applications involving electric heating. However, in applications where heating is an unwanted by-product of current use (e.g., load losses in electrical transformers) the diversion of energy is often referred to as resistive loss. The use of high voltages in electric power transmission systems is specifically designed to reduce such losses in cabling by operating with commensurately lower currents. The ring circuits, or ring mains, used in UK homes are another example, where power is delivered to outlets at lower currents (per wire, by using two paths in parallel), thus reducing Joule heating in the wires. Joule heating does not occur in superconducting materials, as these materials have zero electrical resistance in the superconducting state.
Resistors create electrical noise, called Johnson–Nyquist noise. There is an intimate relationship between Johnson–Nyquist noise and Joule heating, explained by the fluctuation-dissipation theorem.
Formulas
Direct current
The most fundamental formula for Joule heating is the generalized power equation:
where
is the power (energy per unit time) converted from electrical energy to thermal energy,
is the current travelling through the resistor or other element,
is the voltage drop across the element.
The explanation of this formula () is:
Assuming the element behaves as a perfect resistor and that the power is completely converted into heat, the formula can be re-written by substituting Ohm's law, , into the generalized power equation:
where R is the resistance.
Voltage can be increased in DC circuits by connecting batteries or solar panels in series.
Alternating current
When current varies, as it does in AC circuits,
where t is time and P is the instantaneous active power being converted from electrical energy to heat. Far more often, the average power is of more interest than the instantaneous power:
where "avg" denotes average (mean) over one or more cycles, and "rms" denotes root mean square.
These formulas are valid for an ideal resistor, with zero reactance. If the reactance is nonzero, the formulas are modified:
where is phase difference between current and voltage, means real part, Z is the complex impedance, and Y* is the complex conjugate of the admittance (equal to 1/Z*).
For more details in the reactive case, see AC power.
Differential form
Joule heating can also be calculated at a particular location in space. The differential form of the Joule heating equation gives the power per unit volume.
Here, is the current density, and is the electric field. For a material with a conductivity , and therefore
where is the resistivity. This directly resembles the "" term of the macroscopic form.
In the harmonic case, where all field quantities vary with the angular frequency as , complex valued phasors and are usually introduced for the current density and the electric field intensity, respectively. The Joule heating then reads
where denotes the complex conjugate.
Electricity transmission
Overhead power lines transfer electrical energy from electricity producers to consumers. Those power lines have a nonzero resistance and therefore are subject to Joule heating, which causes transmission losses.
The split of power between transmission losses (Joule heating in transmission lines) and load (useful energy delivered to the consumer) can be approximated by a voltage divider. In order to minimize transmission losses, the resistance of the lines has to be as small as possible compared to the load (resistance of consumer appliances). Line resistance is minimized by the use of copper conductors, but the resistance and power supply specifications of consumer appliances are fixed.
Usually, a transformer is placed between the lines and consumption. When a high-voltage, low-intensity current in the primary circuit (before the transformer) is converted into a low-voltage, high-intensity current in the secondary circuit (after the transformer), the equivalent resistance of the secondary circuit becomes higher and transmission losses are reduced in proportion.
During the war of currents, AC installations could use transformers to reduce line losses by Joule heating, at the cost of higher voltage in the transmission lines, compared to DC installations.
Applications
Food processing
Joule heating is a flash pasteurization (also called "high-temperature short-time" (HTST)) aseptic process that runs an alternating current of 50–60 Hz through food. Heat is generated through the food's electrical resistance. As the product heats, electrical conductivity increases linearly. A higher electrical current frequency is best as it reduces oxidation and metallic contamination. This heating method is best for foods that contain particulates suspended in a weak salt-containing medium due to their high resistance properties.
Heat is generated rapidly and uniformly in the liquid matrix as well as in particulates, producing a higher quality sterile product that is suitable for aseptic processing.
Electrical energy is linearly translated to thermal energy as electrical conductivity increases, and this is the key process parameter that affects heating uniformity and heating rate. This heating method is best for foods that contain particulates suspended in a weak salt containing medium due to their high resistance properties. Ohmic heating is beneficial due to its ability to inactivate microorganisms through thermal and non-thermal cellular damage.
This method can also inactivate antinutritional factors thereby maintaining nutritional and sensory properties. However, ohmic heating is limited by viscosity, electrical conductivity, and fouling deposits. Although ohmic heating has not yet been approved by the Food and Drug Administration (FDA) for commercial use, this method has many potential applications, ranging from cooking to fermentation.
There are different configurations for continuous ohmic heating systems, but in the most basic process, a power supply or generator is needed to produce electrical current. Electrodes, in direct contact with food, pass electric current through the matrix. The distance between the electrodes can be adjusted to achieve the optimum electrical field strength.
The generator creates the electrical current which flows to the first electrode and passes through the food product placed in the electrode gap. The food product resists the flow of current causing internal heating. The current continues to flow to the second electrode and back to the power source to close the circuit. The insulator caps around the electrodes controls the environment within the system.
The electrical field strength and the residence time are the key process parameters which affect heat generation.
The ideal foods for ohmic heating are viscous with particulates.
Thick soups
Sauces
Stews
Salsa
Fruit in a syrup medium
Milk
Ice cream mix
Egg
Whey
Heat sensitive liquids
Soymilk
The efficiency by which electricity is converted to heat depends upon on salt, water, and fat content due to their thermal conductivity and resistance factors. In particulate foods, the particles heat up faster than the liquid matrix due to higher resistance to electricity and matching conductivity can contribute to uniform heating. This prevents overheating of the liquid matrix while particles receive sufficient heat processing. Table 1 shows the electrical conductivity values of certain foods to display the effect of composition and salt concentration. The high electrical conductivity values represent a larger number of ionic compounds suspended in the product, which is directly proportional to the rate of heating. This value is increased in the presence of polar compounds, like acids and salts, but decreased with nonpolar compounds, like fats. Electrical conductivity of food materials generally increases with temperature, and can change if there are structural changes caused during heating such as gelatinization of starch. Density, pH, and specific heat of various components in a food matrix can also influence heating rate.
Benefits of Ohmic heating include: uniform and rapid heating (>1°Cs−1), less cooking time, better energy efficiency, lower capital cost, and heating simulataneously throughout food's volume as compared to aseptic processing, canning, and PEF. Volumetric heating allows internal heating instead of transferring heat from a secondary medium. This results in the production of safe, high quality food with minimal changes to structural, nutritional, and organoleptic properties of food. Heat transfer is uniform to reach areas of food that are harder to heat. Less fouling accumulates on the electrodes as compared to other heating methods. Ohmic heating also requires less cleaning and maintenance, resulting in an environmentally cautious heating method.
Microbial inactivation in ohmic heating is achieved by both thermal and non-thermal cellular damage from the electrical field. This method destroys microorganisms due to electroporation of cell membranes, physical membrane rupture, and cell lysis. In electroporation, excessive leakage of ions and intramolecular components results in cell death. In membrane rupture, cells swell due to an increase in moisture diffusion across the cell membrane. Pronounced disruption and decomposition of cell walls and cytoplasmic membranes causes cells to lyse.
Decreased processing times in ohmic heating maintains nutritional and sensory properties of foods. Ohmic heating inactivates antinutritional factors like lipoxigenase (LOX), polyphenoloxidase (PPO), and pectinase due to the removal of active metallic groups in enzymes by the electrical field. Similar to other heating methods, ohmic heating causes gelatinization of starches, melting of fats, and protein agglutination. Water-soluble nutrients are maintained in the suspension liquid allowing for no loss of nutritional value if the liquid is consumed.
Ohmic heating is limited by viscosity, electrical conductivity, and fouling deposits. The density of particles within the suspension liquid can limit the degree of processing. A higher viscosity fluid will provide more resistance to heating, allowing the mixture to heat up quicker than low viscosity products. A food product's electrical conductivity is a function of temperature, frequency, and product composition. This may be increased by adding ionic compounds, or decreased by adding non-polar constituents. Changes in electrical conductivity limit ohmic heating as it is difficult to model the thermal process when temperature increases in multi-component foods.
The potential applications of ohmic heating range from cooking, thawing, blanching, peeling, evaporation, extraction, dehydration, and fermentation. These allow for ohmic heating to pasteurize particulate foods for hot filling, pre-heat products prior to canning, and aseptically process ready-to-eat meals and refrigerated foods. Prospective examples are outlined in Table 2 as this food processing method has not been commercially approved by the FDA. Since there is currently insufficient data on electrical conductivities for solid foods, it is difficult to prove the high quality and safe process design for ohmic heating. Additionally, a successful 12D reduction for C. botulinum prevention has yet to be validated.
Materials synthesis, recovery and processing
Flash joule heating (transient high-temperature electrothermal heating) has been used to synthesize allotropes of carbon, including graphene and diamond. Heating various solid carbon feedstocks (carbon black, coal, coffee grounds, etc.) to temperatures of ~3000 K for 10-150 milliseconds produces turbostratic graphene flakes. FJH has also been used to recover rare-earth elements used in modern electronics from industrial wastes. Beginning from a fluorinated carbon source, fluorinated activated carbon, fluorinated nanodiamond, concentric carbon (carbon shell around a nanodiamond core), and fluorinated flash graphene can be synthesized.
Gallery
Heating efficiency
Heat is not to be confused with internal energy or synonymously thermal energy. While intimately connected to heat, they are distinct physical quantities.
As a heating technology, Joule heating has a coefficient of performance of 1.0, meaning that every joule of electrical energy supplied produces one joule of heat. In contrast, a heat pump can have a coefficient of more than 1.0 since it moves additional thermal energy from the environment to the heated item.
The definition of the efficiency of a heating process requires defining the boundaries of the system to be considered. When heating a building, the overall efficiency is different when considering heating effect per unit of electric energy delivered on the customer's side of the meter, compared to the overall efficiency when also considering the losses in the power plant and transmission of power.
Hydraulic equivalent
In the energy balance of groundwater flow a hydraulic equivalent of Joule's law is used:
where:
= loss of hydraulic energy () due to friction of flow in -direction per unit of time (m/day), comparable to
= flow velocity in -direction (m/day), comparable to
= hydraulic conductivity of the soil (m/day), the hydraulic conductivity is inversely proportional to the hydraulic resistance which compares to
| Physical sciences | Electrical circuits | Physics |
605210 | https://en.wikipedia.org/wiki/Dryopteridaceae | Dryopteridaceae | The Dryopteridaceae are a family of leptosporangiate ferns in the order Polypodiales. They are known colloquially as the wood ferns. In the Pteridophyte Phylogeny Group classification of 2016 (PPG I), the family is placed in the suborder Polypodiineae. Alternatively, it may be treated as the subfamily Dryopteridoideae of a very broadly defined family Polypodiaceae sensu lato.
The family contains about 1700 species and has a cosmopolitan distribution. Species may be terrestrial, epipetric, hemiepiphytic, or epiphytic. Many are cultivated as ornamental plants. The largest genera are Elaphoglossum (600+), Polystichum (260), Dryopteris (225), and Ctenitis (150). These four genera contain about 70% of the species. Dryopteridaceae diverged from the other families in eupolypods I about 100 million years ago.
Description
The rhizomes are often stout, creeping, ascending, or erect, and sometimes scandent or climbing, with nonclathrate scales at apices. Fronds are usually monomorphic, less often dimorphic, or sometimes scaly or glandular, but less commonly hairy. Petioles have numerous round, vascular bundles arranged in a ring, or rarely as few as three; the adaxial bundles are largest. Veins are pinnate or forking, free to variously anastomosing; the areoles occur with or without included veinlets; sori are usually round, acrostichoid (covering the entire abaxial surface of the lamina) in a few lineages; usually indusiate, or sometimes exindusiate. Indusia, when present, are round-reniform or peltate. Sporangia have three-rowed, short to long stalks; spores are reniform, monolete, perine or winged.
Taxonomy
History
In 1990, Karl U. Kramer and coauthors defined the Dryopteridaceae broadly to include the present family, as well as the Woodsiaceae sensu lato, Onocleaceae, and most of Tectariaceae. Molecular phylogenetic studies found Kramer's version of the Dryopteridaceae to be polyphyletic, and it was split up by Smith and others in 2006. The inclusion of Didymochlaena, Hypodematium, and Leucostegia in the Dryopteridaceae is doubtful. If these three are excluded, then the family is strongly supported as monophyletic in cladistic analyses. Some authors have already treated these genera as outside of the Dryopteridaceae.
In 2007, a phylogenetic study of DNA sequences showed that Pleocnemia should be transferred from the Tectariaceae to the Dryopteridaceae. In 2010, in a paper on bolbitidoid ferns, Arthrobotrya was resurrected from Teratophyllum. Later that year, Mickelia was described as a new genus.
Some species have been removed from the genus Oenotrichia because they do not belong there or even in the family Dennstaedtiaceae where Oenotrichia sensu stricto is placed. These species probably belong in the Dryopteridaceae, but have not yet been given a generic name.
In 2012, a phylogenetic study of Dryopteris and its relatives included Acrophorus, Acrorumohra, Diacalpe, Dryopsis, Nothoperanema, and Peranema within that genus. The Flora of China treatment of the family, published in 2013, used phylogenetic results to sink Lithostegia and Phanerophlebiopsis into Arachniodes.
The Dryopteridaceae Herter, under the classification system of Christenhusz and Chase (2014), were submerged as subfamily Dryopteridoideae Link, one of eight subfamilies constituting family Polypodiaceae. This family corresponds to the clade eupolypods I. The Pteridophyte Phylogeny Group classification of 2016 (PPG I) retained the family.
Phylogeny
The following cladogram for the suborder Polypodiineae (eupolypods I), based on the consensus cladogram in the Pteridophyte Phylogeny Group classification of 2016 (PPG I), shows a likely phylogenetic relationship between Dryopteridaceae and the other families of the clade.
Subdivision
The PPG I classification divides the family into three subfamilies, listed below.
Subfamily Polybotryoideae H.M.Liu & X.C.Zhang
Cyclodium C.Presl
Maxonia C.Chr.
Olfersia Raddi
Polybotrya Humb. & Bonpl. ex Willd.
Polystichopsis (J.Sm.) Holttum
Stigmatopteris C.Chr.
Trichoneuron Ching
Subfamily Elaphoglossoideae (Pic.Serm.) Crabbe, Jermy & Mickel
Arthrobotrya J.Sm.
Bolbitis Schott
Elaphoglossum Schott ex J.Sm.
Lastreopsis Ching
Lomagramma J.Sm.
Megalastrum Holttum
Mickelia R.C.Moran, Labiak & Sundue
Parapolystichum (Keyserl.) Ching
Pleocnemia C.Presl
Rumohra Raddi
Teratophyllum Mett. ex Kuhn
Subfamily Dryopteridoideae Link
Arachniodes Blume
Ctenitis (C.Chr.) C.Chr.
Cyrtomium C.Presl
Dryopteris Adans.
Phanerophlebia C.Presl
Polystichum Roth
Didymochlaena has been removed to Didymochlaenaceae, and Hypodematium and Leucostegia to Hypodematiaceae. Aenigmopteris has at times been suggested to belong to this family, on the grounds of its morphological similarity to Ctenitis, but molecular phylogeny has led to its submersion within Tectaria (Tectariaceae). Dryopolystichum has been placed in Lomariopsidaceae.
| Biology and health sciences | Ferns | Plants |
605528 | https://en.wikipedia.org/wiki/Spotted%20dove | Spotted dove | The spotted dove or eastern spotted dove (Spilopelia chinensis) is a small and somewhat long-tailed pigeon that is a common resident breeding bird across its native range on the Indian subcontinent and in East and Southeast Asia. The species has been introduced to many parts of the world and feral populations have become established.
This species was formerly included in the genus Streptopelia with other turtle-doves, but studies suggest that they differ from typical members of that genus. This dove is long tailed buff brown with a white-spotted black collar patch on the back and sides of the neck. The tail tips are white and the wing coverts have light buff spots.
There are considerable plumage variations across populations within its wide range. The species is found in light forests and gardens as well as in urban areas. They fly from the ground with an explosive flutter and will sometimes glide down to a perch. It is also called the mountain dove, pearl-necked dove, lace-necked dove, and spotted turtle-dove.
Taxonomy
The spotted dove was formally described in 1786 by the Austrian naturalist Giovanni Antonio Scopoli and given the binomial name Columba chinensis. Scopoli based his account on "La tourterelle gris de la Chine" that had been described and illustrated in 1782 by the French naturalist Pierre Sonnerat in the second volume of his book Voyage aux Indes orientales et à la Chine. This species was formerly included in the genus Streptopelia. A molecular phylogenetic study published in 2001 found the genus was paraphyletic with respect to Columba. To create monophyletic genera the spotted dove as well as the closely related laughing dove were moved to the resurrected genus Spilopelia that had been introduced by the Swedish zoologist Carl Sundevall in 1873. Sundevall had designated Columba tigrina as the type species, a taxon that is now considered a subspecies of the spotted dove.
Several subspecies have been proposed for the plumage and size variation seen in different geographic populations. The nominate form is from China (Canton), which is also the origin of the introduced population in Hawaii. Subspecies formosa from Taiwan has been considered as doubtful and indistinguishable from the nominate population. The population in India suratensis (type locality Surat) and ceylonensis from Sri Lanka have fine rufous or buff spots on the back. There is a size reduction trend with specimens from southern India being smaller, and ceylonensis may merely be a part of this cline. The lesser and median wing-coverts are also spotted at the tip in buff. This spotting is lacking on populations further north and east of India, such as tigrina, which also differ greatly in vocalizations from the Indian forms. The population from Hainan Island is placed in hainana. Others like vacillans (=chinensis) and forresti (= tigrina) and edwardi (from Chabua = suratensis) have been considered invalid.
Five subspecies are recognised:
Spilopelia chinensis suratensis (Gmelin, JF, 1789) – Pakistan, India, Nepal and Bhutan
Spilopelia chinensis ceylonensis (Reichenbach, 1851) – Sri Lanka (has shorter wings than suratensis)
Spilopelia chinensis tigrina (Temminck, 1809) – Bangladesh and northeast India through Indochina to Philippines and the Sunda Islands
Spilopelia chinensis chinensis (Scopoli, 1786) – northeast Myanmar to central and east China, Taiwan
Spilopelia chinensis hainana (Hartert, 1910) – Hainan (off southeast China)
The subspecies S. c. suratensis and S. c. ceylonensis differ significantly from the other subspecies in both plumage and vocalization. This has led some ornithologists to treat S. c. suratensis as a separate species, the western spotted dove.
Description
The ground colour of this long and slim dove is rosy buff below shading into grey on the head and belly. There is a half collar on the back and sides of the neck made of black feathers that bifurcate and have white spots at the two tips. The median coverts have brown feathers tipped with rufous spots in the Indian and Sri Lankan subspecies which are divided at the tip by a widening grey shaft streak.
The wing feathers are dark brown with grey edges. The centre of the abdomen and vent are white. The outer tail feathers are tipped in white and become visible when the bird takes off. Sexes are similar, but juveniles are duller than adults and do not acquire the neck spots until they are mature. The length ranges from 28 to 32 centimetres (11.2 to 12.8 inches).
Abnormal plumages such as leucism can sometimes occur in the wild.
Distribution and habitat
The spotted dove in its native range in Asia is found across a range of habitats including woodland, scrub, farmland and habitation. In India it tends to be found in the moister regions, with the laughing dove (S. senegalensis) appearing more frequently in drier areas. These doves are mostly found on the ground where they forage for seeds and grain or on low vegetation.
The species has become established in many areas outside its native range. These areas include Hawaii, southern California, Mauritius, Australia and New Zealand.
In Australia they were introduced into Melbourne in the 1860s and have since spread but there is insufficient evidence that they compete with native doves. They are now found in streets, parks, gardens, agricultural areas, and tropical scrubs in diverse locations throughout eastern Australia and around the cities and major towns across southern Australia. The original populations appear to be S. c. chinensis and S. c. tigrina in varying proportions.
Behaviour and ecology
Spotted doves move around in pairs or small groups as they forage on the ground for grass seeds, grains, fallen fruits and seeds of other plants. They may however take insects occasionally and have been recorded feeding on winged termites. The flight is quick with regular beats and an occasional sharp flick of the wings. A display flight involves taking off at a steep angle with a loud clapping of the wing and then slowly gliding down with the tail spread out.
The breeding season is spread out in warm regions but tends to be in summer in the temperate ranges. In Hawaii, they breed all year round, as do all three other introduced species of doves. Males coo, bow and make aerial displays in courtship. In southern Australia, they breed mostly from September to January, and in the north in autumn. They nest mainly in low vegetation, building a flimsy cup of twigs in which two whitish eggs are laid. Nests are sometimes placed on the ground or on buildings and other structures. Both parents take part in building the nest, incubating and feeding the young. The eggs hatch after about 13 days and fledge after a fortnight. More than one brood may be raised.
The vocalizations of the spotted dove include cooing softly with a Krookruk-krukroo... kroo kroo kroo with the number of terminal kroos varying in the Indian population and absent in tigrina, chinensis and other populations to the east.
The species has been extending its range in many parts of the world. Populations may sometimes rise and fall rapidly, within a span of about five years. In the Philippines, the species may be outcompeting the Streptopelia dusumieri. Their habit of flushing into the air when disturbed makes them a hazard on airfields, often colliding with aircraft and sometimes causing damage.
| Biology and health sciences | Columbimorphae | Animals |
605587 | https://en.wikipedia.org/wiki/Laughing%20dove | Laughing dove | The laughing dove (Spilopelia senegalensis) is a small pigeon that is a resident breeder in Africa, the Middle East, South Asia, and Western Australia where it has established itself in the wild after being released from Perth Zoo in 1898. This small long-tailed dove is found in dry scrub and semi-desert habitats where pairs can often be seen feeding on the ground. It is closely related to the spotted dove (Spilopelia chinensis) which is distinguished by a white and black chequered necklace. Other names include laughing turtle dove, palm dove and Senegal dove while in Asia the name little brown dove is often used.
Taxonomy
In 1760 the French zoologist Mathurin Jacques Brisson included a description of the laughing dove in his six volume Ornithologie based on a specimen collected in Senegal. He used the French name La tourterelle à gorge tachetée du Sénégal and the Latin Tutur gutture maculato senegalensis. Although Brisson coined Latin names, these do not conform to the binomial system and are not recognised by the International Commission on Zoological Nomenclature. When in 1766 the Swedish naturalist Carl Linnaeus updated his Systema Naturae for the twelfth edition, he added 240 species that had been previously described by Brisson. One of these was the laughing dove which he placed with all the other pigeons in the genus Columba. Linnaeus included a brief description, coined the binomial name Columba senegalensis and cited Brisson's work.
For many years the laughing dove was placed in the genus Streptopelia. A molecular phylogenetic study published in 2001 found this genus was paraphyletic with respect to Columba. To create monophyletic genera the laughing dove as well as the closely related spotted dove were moved to the resurrected genus Spilopelia that had been introduced by the Swedish zoologist Carl Sundevall in 1873.
Five populations with small plumage and size differences have been given the status of subspecies:
S. s. phoenicophila (Hartert, 1916) – Morocco to northwest Libya
S. s. aegyptiaca (Latham, 1790) – Nile Valley (Egypt)
S. s. senegalensis (Linnaeus, 1766) – southern laughing dove, west Arabia, Socotra Island, Africa south of the Sahara
S. s. cambayensis (Gmelin, JF, 1789) – east Arabia and east Iran to Pakistan, India and Bangladesh
S. s. ermanni (Bonaparte, 1856) – Kazakhstan, north Afghanistan, west China
Several other subspecies have been described but are not now generally recognised. These include S. s. sokotrae on Socotra Island, S. s. dakhlae in the Dakhla Oasis, Egypt and S. s. thome on São Tomé Island.
Description
The laughing dove is a long-tailed, slim pigeon, typically in length. It is pinkish brown on the underside with a lilac tinged head and neck. The head and underparts are pinkish, shading to buff on the lower abdomen. A chequered rufous and grey patch is found on the sides of the neck of adults and is made up of split feathers. The upper parts are brownish with a bluish-grey band along the wing. The back is uniform and dull brown in the South Asian population. The African populations S. s. senegalensis and S. s. phoenicophila have a bluish grey rump and upper tail coverts but differ in the shades of the neck and wing feathers while S. s. aegyptiaca is larger and the head and nape are vinous and upper wing coverts are rufous. The tail is graduated and the outer feathers are tipped in white. The sexes are indistinguishable in the field. Young birds lack the chequered neck markings. The legs are red. The populations vary slightly in plumage with those from more arid zones being paler. Abnormal leucistic plumages have been noted.
The chuckling call is a low rolling croo-doo-doo-doo-doo with a rising and falling amplitude.
Distribution and habitat
It is a common and widespread species in scrub, dry farmland, and in areas of human habitation, often becoming very tame. Its range includes much of Sub-Saharan Africa, Saudi Arabia, Iran, Iraq, Afghanistan, Pakistan, and India. It is also found in Cyprus, Palestine, Israel, Lebanon, Syria, Jordan, the UAE, and Turkey (these populations may be derived from human introductions). They are mostly sedentary but some populations may make movements. Birds ringed in Gujarat have been recovered 200 km north in Pakistan and exhausted birds have been recorded landing on ships in the Arabian Sea. The species (thought to belong to the nominate population) was introduced to Perth in 1889 and has become established around Western Australia. Birds that land on ships may be introduced to new regions.
Behaviour and ecology
The species is usually seen in pairs or small parties and only rarely in larger groups. Larger groups are formed especially when drinking at waterholes in arid regions. Small numbers assemble on trees near waterholes before flying to the water's edge where they are able to suck up water like other members of the pigeon family. Laughing doves eat the fallen seeds, mainly of grasses, other vegetable matter and small ground insects such as termites and beetles. They are fairly terrestrial, foraging on the ground in grasslands and cultivation. Their flight is quick and direct with the regular beats and an occasional sharp flick of the wings characteristic of pigeons in general.
The male in courtship display follows the female with head bobbing displays while cooing. The male pecks its folded wings in "displacement-preening" to solicit copulation from the female. A female accepts by crouching and begging for food. The male may indulge in courtship feeding before mounting and copulating. Pairs may preen each other. Males may also launch into the air with wing clapping above their backs and then glide down in a gentle arc when displaying. The species has a spread out breeding season in Africa. Almost year-round in Malawi and Turkey; and mainly May to November in Zimbabwe, February to June in Egypt and Tunisia. In Australia the main breeding season is September to November. The nest is a very flimsy platform of twigs built in a low bush and sometimes in crevices or under the eaves of houses. Both parents build the nest with males bringing the twigs which are then placed by the female. Two eggs are laid within an interval of a day between them and both parents take part in building the nest, incubating and feeding the young. Males spend more time incubating the nest during the day. The eggs are incubated after the second egg is laid and the eggs hatch after about 13 to 15 days. Nesting adults may feign injury to distract and draw predators away from the nest. Multiple broods may be raised by the same pair in the same nest. Seven broods by the same pair have been noted in Turkey. Initially the altricial hatchlings are fed with regurgitated crop-milk, a secretion from the lining of the crop of parent birds. The young fledge and leave the nest after about 14 to 16 days. The Jacobin cuckoo sometimes lays its egg in the nests of the laughing dove in Africa.
Feral populations in Australia are sometimes infected by a virus that causes symptoms similar to that produced in parrots by psittacine beak and feather disease. Several ectoparasitic bird lice have been found on the species and include those in the genera Coloceras, Columbicola, Bonomiella and Hohorstiella. A blood parasite Trypanosoma hannae has been recorded in the species. Southern grey shrike have been observed preying on an adult laughing dove in northwestern India while the lizard buzzard is a predator of the species in Africa. South African birds sometimes show a beak deformity in which the upper mandible overgrowth occurs.
| Biology and health sciences | Columbimorphae | Animals |
606149 | https://en.wikipedia.org/wiki/Meso%20compound | Meso compound | A meso compound or meso isomer is an optically inactive isomer in a set of stereoisomers, at least two of which are optically active. This means that despite containing two or more stereocenters, the molecule is not chiral. A meso compound is superposable on its mirror image (not to be confused with superimposable, as any two objects can be superimposed over one another regardless of whether they are the same). Two objects can be superposed if all aspects of the objects coincide and it does not produce a "(+)" or "(-)" reading when analyzed with a polarimeter. The name is derived from the Greek mésos meaning “middle”.
For example, tartaric acid can exist as any of three stereoisomers depicted below in a Fischer projection. Of the four colored pictures at the top of the diagram, the first two represent the meso compound (the 2R,3S and 2S,3R isomers are equivalent), followed by the optically active pair of levotartaric acid (L-(R,R)-(+)-tartaric acid) and dextrotartaric acid (D-(S,S)-(-)-tartaric acid). The meso compound is bisected by an internal plane of symmetry that is not present for the non-meso isomers (indicated by an X). That is, on reflecting the meso compound through a mirror plane perpendicular to the screen, the same stereochemistry is obtained; this is not the case for the non-meso tartaric acid, which generates the other enantiomer. The meso compound must not be confused with a 50:50 racemic mixture of the two optically-active compounds, although neither will rotate light in a polarimeter.
It is a requirement for two of the stereocenters in a meso compound to have at least two substituents in common (although having this characteristic does not necessarily mean that the compound is meso). For example, in 2,4-pentanediol, both the second and fourth carbon atoms, which are stereocenters, have all four substituents in common.
Since a meso isomer has a superposable mirror image, a compound with a total of n chiral centers cannot attain the theoretical maximum of 2n stereoisomers if one of the stereoisomers is meso.
A meso isomer need not have a mirror plane. It may have an inversion or a rotoreflexion symmetry such as S. For example, there are two meso isomers of 1,4-difluoro-2,5-dichlorocyclohexane but neither has a mirror plane, and there are two meso isomers of 1,2,3,4-tetrafluorospiropentane (see figure).
Cyclic meso compounds
1,2-substituted cyclopropane has a meso cis-isomer (molecule has a mirror plane) and two trans-enantiomers:
The two cis stereoisomers of 1,2-substituted cyclohexanes behave like meso compounds at room temperature in most cases. At room temperature, most 1,2-disubstituted cyclohexanes undergo rapid ring flipping (exceptions being rings with bulky substituents), and as a result, the two cis stereoisomers behave chemically identically with chiral reagents. At low temperatures, however, this is not the case, as the activation energy for the ring-flip cannot be overcome, and they therefore behave like enantiomers. Also noteworthy is the fact that when a cyclohexane undergoes a ring flip, the absolute configurations of the stereocenters do not change.
| Physical sciences | Stereochemistry | Chemistry |
1018642 | https://en.wikipedia.org/wiki/Seawall | Seawall | A seawall (or sea wall) is a form of coastal defense constructed where the sea, and associated coastal processes, impact directly upon the landforms of the coast. The purpose of a seawall is to protect areas of human habitation, conservation, and leisure activities from the action of tides, waves, or tsunamis. As a seawall is a static feature, it will conflict with the dynamic nature of the coast and impede the exchange of sediment between land and sea.
Seawall designs factor in local climate, coastal position, wave regime (determined by wave characteristics and effectors), and value (morphological characteristics) of landform. Seawalls are hard engineering shore-based structures that protect the coast from erosion. Various environmental issues may arise from the construction of a seawall, including the disruption of sediment movement and transport patterns. Combined with a high construction cost, this has led to increasing use of other soft engineering coastal management options such as beach replenishment.
Seawalls are constructed from various materials, most commonly reinforced concrete, boulders, steel, or gabions. Other possible construction materials include vinyl, wood, aluminum, fiberglass composite, and biodegradable sandbags made of jute and coir. In the UK, seawall also refers to an earthen bank used to create a polder, or a dike construction. The type of material used for construction is hypothesized to affect the settlement of coastal organisms, although the precise mechanism has yet to be identified.
Types
A seawall works by reflecting incident wave energy back into the sea, thus reducing the energy available to cause erosion. Seawalls have two specific weaknesses. Wave reflection from the wall may result in hydrodynamic scour and subsequent lowering of the sand level of the fronting beach. Seawalls may also accelerate the erosion of adjacent, unprotected coastal areas by affecting the littoral drift process.
Different designs of man-made tsunami barriers include building reefs and forests to above-ground and submerged seawalls. Starting just weeks after the disaster, in January 2005, India began planting Casuarina and coconut saplings on its coast as a natural barrier against future disasters like the 2004 Indian Ocean earthquake. Studies have found that an offshore tsunami wall could reduce tsunami wave heights by up to 83%.
The appropriate seawall design relies on location-specific aspects, including surrounding erosion processes. There are three main types of seawalls: vertical, curved, stepped, and mounds (see table below).
Natural barriers
A report published by the United Nations Environment Programme (UNEP) suggests that the tsunami of 26 December 2004 caused less damage in the areas where natural barriers were present, such as mangroves, coral reefs or coastal vegetation. A Japanese study of this tsunami in Sri Lanka used satellite imagery modelling to establish the parameters of coastal resistance as a function of different types of trees. Natural barriers, such as coral reefs and mangrove forests, prevent the spread of tsunamis and the flow of coastal waters and mitigated the flood and surge of water.
Trade-offs
A cost-benefit approach is an effective way to determine whether a seawall is appropriate and whether the benefits are worth the expense. Besides controlling erosion, consideration must be given to the effects of hardening a shoreline on natural coastal ecosystems and human property or activities. A seawall is a static feature which can conflict with the dynamic nature of the coast and impede the exchange of sediment between land and sea. The table below summarizes some positive and negative effects of seawalls which can be used when comparing their effectiveness with other coastal management options, such as beach nourishment.
Generally, seawalls can be a successful way to control coastal erosion, but only if they are constructed well and out of materials that can withstand the force of ongoing wave energy. Some understanding is needed of the coastal processes and morphodynamics specific to the seawall location. Seawalls can be very helpful; they can offer a more long-term solution than soft engineering options, additionally providing recreation opportunities and protection from extreme events as well as everyday erosion. Extreme natural events expose weaknesses in the performance of seawalls, and analyses of these can lead to future improvements and reassessment.
Issues
Sea level rise
Sea level rise creates an issue for seawalls worldwide as it raises both the mean normal water level and the height of waves during extreme weather events, which the current seawall heights may be unable to cope with. The most recent analyses of long, good-quality tide gauge records (corrected for GIA and when possible for other vertical land motions by the Global Positioning System, GPS) indicate a mean rate of sea level rise of 1.6–1.8 mm/yr over the twentieth century. The Intergovernmental Panel on Climate Change (IPCC) (1997) suggested that sea level rise over the next 50 – 100 years will accelerate with a projected increase in global mean sea level of +18 cm by 2050 AD. This data is reinforced by Hannah (1990) who calculated similar statistics including a rise of between +16-19.3 cm throughout 1900–1988. Superstorm Sandy of 2012 is an example of the devastating effects rising sea levels can cause when mixed with a perfect storm. Superstorm Sandy sent a storm surge of 4–5 m onto New Jersey's and New York's barrier island and urban shorelines, estimated at $70 billion in damage. This problem could be overcome by further modeling and determining the extension of height and reinforcement of current seawalls which needs to occur for safety to be ensured in both situations. Sea level rise also will cause a higher risk of flooding and taller tsunamis.
Hydrostatic water pressure
Seawalls, like all retaining walls, must relieve the buildup of water pressure. Water pressure buildup is caused when groundwater is not drained from behind the seawall. Groundwater against a seawall can be from the area's natural water-table, rain percolating into the ground behind the wall and waves overtopping the wall. The water table can also rise during periods of high water (high tide). Lack of adequate drainage can cause the seawall to buckle, move, bow, crack, or collapse. Sinkholes may also develop as the escaping water pressure erodes soil through or around the drainage system.
Extreme events
Extreme events also pose a problem as it is not easy for people to predict or imagine the strength of hurricane or storm-induced waves compared to normal, expected wave patterns. An extreme event can dissipate hundreds of times more energy than everyday waves, and calculating structures that will stand the force of coastal storms is difficult and, often the outcome can become unaffordable. For example, the Omaha Beach seawall in New Zealand was designed to prevent erosion from everyday waves only, and when a storm in 1976 carved out ten meters behind the existing seawall, the whole structure was destroyed.
Ecosystem impacts
The addition of seawalls near marine ecosystems can lead to increased shadowing effects in the waters surrounding the seawall. Shadowing reduces the light and visibility within the water, which may disrupt the distribution as well as foraging capabilities of certain species. The sediment surrounding seawalls tends to have less favorable physical properties (Higher calcification levels, less structural organization of crystalline structure, low silicon content, and less macroscale roughness) when compared to natural shorelines, which can present issues for species that reside on the seafloor.
The Living Seawalls project, which was launched in Sydney, Australia, in 2018, aims to help many of the marine species in Sydney Harbour to flourish, thus enhancing its biodiversity, by modifying the design of its seawalls. It entails covering parts of the seawalls with specially-designed tiles that mimic natural microhabitats - with crevices and other features that more closely resemble natural rocks. In September 2021, the Living Seawalls project was announced as a finalist for the international environment award the Earthshot Prize. Since 2022 it has become part of Project Restore, under the auspices of the Sydney Institute of Marine Science.
Other issues
Some further issues include a lack of long-term trend data of seawall effects due to a relatively short duration of data records; modeling limitations and comparisons of different projects and their effects being invalid or unequal due to different beach types; materials; currents; and environments. Lack of maintenance is also a major issue with seawalls. In 2013, more than 5,000 feet (1,500 m) of seawall was found to be crumbling in Punta Gorda, Florida. Residents of the area pay hundreds of dollars each year for a seawall repair program. The problem is that most of the seawalls are over a half-century old and are being destroyed by only heavy downpours. If not kept in check, seawalls lose effectiveness and become expensive to repair.
History and examples
Seawall construction has existed since ancient times. In the first century BCE, Romans built a seawall or breakwater at Caesarea Maritima
creating an artificial harbor (Sebastos Harbor). The construction used Pozzolana concrete which hardens in contact with seawater. Barges were constructed and filled with the concrete. They were floated into position and sunk. The resulting harbor/breakwater/seawall is still in existence today – more than 2000 years later.
The oldest known coastal defense is believed to be a 100-meter row of boulders in the Mediterranean Sea off the coast of Israel. Boulders were positioned in an attempt to protect the coastal settlement of Tel Hreiz from sea rise following the last glacial maximum. Tel Hreiz was discovered in 1960 by divers searching for shipwrecks, but the row of boulders was not found until storms cleared a sand cover in 2012.
More recently, seawalls were constructed in 1623 in Canvey Island, UK, when great floods of the Thames estuary occurred, prompting the construction of protection for further events in this flood-prone area. Since then, seawall design has become more complex and intricate in response to an improvement in materials, technology, and an understanding of how coastal processes operate. This section will outline some key case studies of seawalls in chronological order and describe how they have performed in response to tsunamis or ongoing natural processes and how effective they were in these situations. Analyzing the successes and shortcomings of seawalls during severe natural events allows their weaknesses to be exposed, and areas become visible for future improvement.
Canada
The Vancouver Seawall is a stone seawall constructed around the perimeter of Stanley Park in Vancouver, British Columbia. The seawall was constructed initially as waves created by ships passing through the First Narrows eroding the area between Prospect Point and Brockton Point. Construction of the seawall began in 1917, and since then this pathway has become one of the most used features of the park by both locals and tourists and now extends 22 km in total. The construction of the seawall also provided employment for relief workers during the Great Depression and seamen from on Deadman's Island who were facing punishment detail in the 1950s (Steele, 1985).
Overall, the Vancouver Seawall is a prime example of how seawalls can simultaneously provide shoreline protection and a source of recreation which enhances human enjoyment of the coastal environment. It also illustrates that although shoreline erosion is a natural process, human activities, interactions with the coast, and poorly planned shoreline development projects can accelerate natural erosion rates.
India
On December 26, 2004, towering waves of the 2004 Indian Ocean earthquake tsunami crashed against India's south-eastern coastline killing thousands. However, the former French colonial enclave of Pondicherry escaped unscathed. This was primarily due to French engineers who had constructed (and maintained) a massive stone seawall during the time when the city was a French colony. This 300-year-old seawall effectively kept Pondicherry's historic center dry even though tsunami waves drove water above the normal high-tide mark.
The barrier was initially completed in 1735 and over the years, the French continued to fortify the wall, piling huge boulders along its coastline to stop erosion from the waves pounding the harbor. At its highest, the barrier running along the water's edge reaches about above sea level. The boulders, some weighing up to a ton, are weathered black and brown. The seawall is inspected every year and whenever gaps appear or the stones sink into the sand, the government adds more boulders to keep it strong.
The Union Territory of Pondicherry recorded around 600 deaths from the huge tsunami waves that struck India's coast after the mammoth underwater earthquake (which measured 9.0 on the moment magnitude scale) off Indonesia, but most of those killed were fishermen who lived in villages beyond the artificial barrier which reinforces the effectiveness of seawalls.
Japan
At least 43 percent of Japan's coastline is lined with concrete seawalls or other structures designed to protect the country against high waves, typhoons, or even tsunamis. During the 2011 Tōhoku earthquake and tsunami, the seawalls in most areas were overwhelmed. In Kamaishi, waves surmounted the seawall – the world's largest, erected a few years ago in the city's harbor at a depth of , a length of and a cost of $1.5 billion – and eventually submerged the city center.
The risks of dependence on seawalls were most evident in the crisis at the Fukushima Dai-ichi and Fukushima Dai-ni nuclear power plants, both located along the coast close to the earthquake zone, as the tsunami washed over walls that were supposed to protect the plants. Arguably, the additional defense provided by the seawalls presented an extra margin of time for citizens to evacuate and also stopped some of the full force of energy which would have caused the wave to climb higher in the backs of coastal valleys. In contrast, the seawalls also acted in a negative way to trap water and delay its retreat.
The failure of the world's largest seawall, which cost $1.5 billion to construct, shows that building stronger seawalls to protect larger areas would have been even less cost-effective. In the case of the ongoing crisis at the nuclear power plants, higher and stronger seawalls should have been built if power plants were to be built at that site. Fundamentally, the devastation in coastal areas and a final death toll predicted to exceed 10,000 could push Japan to redesign its seawalls or consider more effective alternative methods of coastal protection for extreme events. Such hardened coastlines can also provide a false sense of security to property owners and local residents as evident in this situation.
Seawalls along the Japanese coast have also been criticized for cutting settlements off from the sea, making beaches unusable, presenting an eyesore, disturbing wildlife, and being unnecessary.
United States
After 2012's Hurricane Sandy, New York City Mayor Bill de Blasio invested $3,000,000,000 in a hurricane restoration fund, with part of the money dedicated to building new seawalls and protection from future hurricanes.
A New York Harbor Storm-Surge Barrier has been proposed, but not voted on or funded by Congress or the State of New York.
In Florida, tiger dams are used to protect homes near the coast.
| Technology | Coastal infrastructure | null |
1018849 | https://en.wikipedia.org/wiki/Perseus%20Arm | Perseus Arm | The Perseus Arm is one of two major spiral arms of the Milky Way galaxy. The second major arm is called the Scutum–Centaurus Arm. The Perseus Arm begins from the distal end of the long Milky Way central bar. Previously thought to be 13,000 light-years away, it is now thought to lie 6,400 light years from the Solar System.
Overview
The Milky Way is a barred spiral galaxy with two major arms and a number of minor arms or spurs. The Perseus Spiral Arm, with a radius of approximately 10.7 kiloparsecs, is located between the minor Cygnus and Carina–Sagittarius Arms. It is named after the Perseus constellation in the direction of which it is seen from Earth.
Recently, scientists in two large radio astronomy projects, the Bar and Spiral Structure Legacy (BeSSeL) Survey and the Japanese VLBI Exploration of Radio Astrometry (VERA), have made great efforts over about 20 years to measure the trigonometric parallaxes toward about 200 water vapor () and methanol () masers in massive star-forming regions in the Milky Way. They have employed these parallax measurements to delineate the forms of spiral arms from the Galactic longitude 2 to 240 degrees and extended the spiral arm traces into the portion of the Milky Way seen from the Southern Hemisphere using tangencies along some arms based on carbon monoxide emission. The image clearly presents the Milky Way as a barred spiral galaxy with fairly symmetric four major arms and some extra arm segments and spurs.
The Perseus Arm is one of the four major arms. The arm is the length of more than 60,000 lr and the width of about 1,000 lr and the spiral extension in the pitch angle near 9 degree.
There is speculation that the local spur known as the Orion–Cygnus Arm, which includes the Solar System and Earth and is located inside of the Perseus Arm, or is a branch of it, but this is unconfirmed.
The Perseus Arm contains the Double Cluster and a number of Messier objects:
The Crab Nebula (M1)
Open Cluster M36
Open Cluster M37
Open Cluster M38
Open Cluster M52
Open Cluster M103.
| Physical sciences | Milky Way | Astronomy |
1018854 | https://en.wikipedia.org/wiki/Scutum%E2%80%93Centaurus%20Arm | Scutum–Centaurus Arm | The Scutum–Centaurus Arm, also known as Scutum-Crux arm, is a long, diffuse curving streamer of stars, gas and dust that spirals outward from the proximate end of the Milky Way's central bar. The Milky Way has been posited since the 1950s to have four spiral arms; numerous studies contest or nuance this number. In 2008, observations using the Spitzer Space Telescope failed to show the expected density of red clump giants in the direction of the Sagittarius and Norma arms. In January 2014, a 12-year study into the distribution and lifespan of massive stars and a 2013-reporting study of the distribution of masers and open clusters both found corroboratory, though would not state irrefutable, evidence for four principal spiral arms.
The Scutum–Centaurus Arm lies between the minor Carina–Sagittarius Arm and the minor Norma Arm. The Scutum–Centaurus Arm starts near the core as the Scutum Arm, then gradually turns into the Centaurus Arm.
The region where the Scutum–Centaurus Arm connects to the bar of the galaxy is rich in star-forming regions and open clusters. In 2006 a large cluster of new stars containing 14 red supergiant stars was discovered there and named RSGC1. In 2007 a cluster of approximately 50,000 newly formed stars named RSGC2 was located only a few hundred light years from RSGC1. It is thought to be less than 20 million years old and contains 26 red supergiant stars, the largest grouping of such stars known. Other clusters in this region include RSGC3 and Alicante 8.
| Physical sciences | Milky Way | Astronomy |
1019118 | https://en.wikipedia.org/wiki/Crown%20group | Crown group | In phylogenetics, the crown group or crown assemblage is a collection of species composed of the living representatives of the collection, the most recent common ancestor of the collection, and all descendants of the most recent common ancestor. It is thus a way of defining a clade, a group consisting of a species and all its extant or extinct descendants. For example, Neornithes (birds) can be defined as a crown group, which includes the most recent common ancestor of all modern birds, and all of its extant or extinct descendants.
The concept was developed by Willi Hennig, the formulator of phylogenetic systematics, as a way of classifying living organisms relative to their extinct relatives in his "Die Stammesgeschichte der Insekten",
and the "crown" and "stem" group terminology was coined by R. P. S. Jefferies in 1979. Though formulated in the 1970s, the term was not commonly used until its reintroduction in 2000 by Graham Budd and Sören Jensen.
Contents of the crown group
It is not necessary for a species to have living descendants in order for it to be included in the crown group. Extinct side branches on the family tree that are descended from the most recent common ancestor of living members will still be part of a crown group. For example, if we consider the crown-birds (i.e. all extant birds and the rest of the family tree back to their most recent common ancestor), extinct side branches like the dodo or great auk are still descended from the most recent common ancestor of all living birds, so fall within the bird crown group. One very simplified cladogram for birds is shown below:
In this diagram, the clade labelled "Neornithes" is the crown group of birds: it includes the most recent common ancestor of all living birds and its descendants, living or not. Although considered to be birds (i.e. members of the clade Aves), Archaeopteryx and other extinct groups are not included in the crown group, as they fall outside the Neornithes clade, being descended from an earlier ancestor.
An alternative definition does not require any members of a crown group to be extant, only to have resulted from a "major cladogenesis event". The first definition forms the basis of this article.
Often, the crown group is given the designation "crown-", to separate it from the group as commonly defined. Both birds and mammals are traditionally defined by their traits, and contain fossil members that lived before the last common ancestors of the living groups or, like the mammal Haldanodon, were not descended from that ancestor although they lived later. Crown-Aves and Crown-Mammalia therefore differ slightly in content from the common definition of Aves and Mammalia. This has caused some confusion in the literature.
Other groups under the crown group concept
The cladistic idea of strictly using the topology of the phylogenetic tree to define groups necessitates other definitions than crown groups to adequately define commonly discussed fossil groups. Thus, a host of prefixes have been defined to describe various branches of the phylogenetic tree relative to extant organisms.
Pan-group
A pan-group or total group is the crown group and all organisms more closely related to it than to any other extant organisms. In a tree analogy, it is the crown group and all branches back to (but not including) the split with the closest branch to have living members. The Pan-Aves thus contain the living birds and all (fossil) organisms more closely related to birds than to crocodilians (their closest living relatives). The phylogenetic lineage leading back from Neornithes to the point where it merges with the crocodilian lineage, along with all side branches, constitutes pan-birds. In addition to non-crown group primitive birds like Archaeopteryx, Hesperornis and Confuciusornis, therefore, pan-group birds would include all dinosaurs and pterosaurs as well as an assortment of non-crocodilian animals like Marasuchus.
Pan-Mammalia consists of all mammals and their fossil ancestors back to the phylogenetic split from the remaining amniotes (the Sauropsida). Pan-Mammalia is thus an alternative name for Synapsida.
Stem groups
A stem group is a paraphyletic assemblage composed of the members of a pan-group or total group, above, minus the crown group itself (and therefore minus all living members of the pan-group). This leaves primitive relatives of the crown groups, back along the phylogenetic line to (but not including) the last common ancestor of the crown group and their closest living relatives. It follows from the definition that all members of a stem group are extinct. The "stem group" is the most used and most important of the concepts linked to crown groups, as it offers a means to reify and name paraphyletic assemblages of fossils that otherwise do not fit into systematics based on living organisms.
While often attributed to Jefferies (1979), Willmann (2003) traced the origin of the stem group concept to Austrian systematist Othenio Abel (1914), and it was discussed and diagrammed in English as early as 1933 by A. S. Romer.
Alternatively, the term "stem group" is sometimes used in a wider sense to cover any members of the traditional taxon falling outside the crown group. Permian synapsids like Dimetrodon or Anteosaurus are stem mammals in the wider sense but not in the narrower one.
Often, an (extinct) grouping is identified as belonging together. Later, it may be realized other (extant) groupings actually emerged within such grouping, rendering them a stem grouping. Cladistically, the new groups should then be added to the group, as paraphyletic groupings are not natural. In any case, stem groupings with living descendants should not be viewed as a cohesive group, but their tree should be further resolved to reveal the full bifurcating phylogeny.
Examples of stem groups (in the wider sense)
Stem birds perhaps constitute the most cited example of a stem group, as the phylogeny of this group is fairly well known. The following cladogram, based on Benton (2005), illustrates the concept:
The crown group here is Neornithes, all modern bird lineages back to their last common ancestor. The closest living relatives of birds are crocodilians. If we follow the phylogenetic lineage leading to Neornithes to the left, the line itself and all side branches belong to the stem birds until the lineage merges with that of the crocodilians. In addition to non-crown group primitive birds like Archaeopteryx, Hesperornis and Confuciusornis, stem group birds include the dinosaurs and the pterosaurs. The last common ancestor of birds and crocodilians—the first crown group archosaur—was neither bird nor crocodilian and possessed none of the features unique to either. As the bird stem group evolved, distinctive bird features such as feathers and hollow bones appeared. Finally, at the base of the crown group, all traits common to extant birds were present.
Under the widely used total-group perspective, the Crocodylomorpha would become synonymous with the Crocodilia, and the Avemetatarsalia would become synonymous with the birds, and the above tree could be summarized as
An advantage of this approach is that declaring Theropoda to be birds (or Pan-aves) is more specific than declaring it to be a member of the Archosauria, which would not exclude it from the Crocodilia branch. Basal branch names such as Avemetatarsalia are usually more obscure. However, not so advantageous are the facts that "Pan-Aves" and "Aves" are not the same group, the circumscription of the concept of "Pan-Aves" (synonymous with Avemetatarsalia) is only evident by examination of the above tree, and calling both groups "birds" is ambiguous.
Stem mammals are those in the lineage leading to living mammals, together with side branches, from the divergence of the lineage from the Sauropsida to the last common ancestor of the living mammals. This group includes the synapsids as well as mammaliaforms like the morganucodonts and the docodonts; the latter groups have traditionally and anatomically been considered mammals even though they fall outside the crown group mammals.
Stem tetrapods are the animals belonging to the lineage leading to tetrapods from their divergence from the lungfish, our nearest relatives among the fishes. In addition to a series of lobe-finned fishes, they also include some of the early labyrinthodonts. Exactly what labyrinthodonts are in the stem group tetrapods rather than the corresponding crown group is uncertain, as the phylogeny of early tetrapods is not well understood. This example shows that crown and stem group definitions are of limited value when there is no consensus phylogeny.
Stem arthropods constitute a group that has seen attention in connection with the Burgess Shale fauna. Several of the finds, including the enigmatic Opabinia and Anomalocaris have some, though not all, features associated with arthropods, and are thus considered stem arthropods. The sorting of the Burgess Shale fauna into various stem groups finally enabled phylogenetic sorting of this enigmatic assemblage and also allowed for identifying velvet worms as the closest living relatives of arthropods.
Stem priapulids are other early Cambrian to middle Cambrian faunas, appearing in Chengjiang to Burgess Shale. The genus Ottoia has more or less the same build as modern priapulids, but phylogenetic analysis indicates that it falls outside the crown group, making it a stem priapulid.
Plesion-group
The name plesion has a long history in biological systematics, and plesion group has acquired several meanings over the years. One use is as "nearby group" (plesion means close to in Greek), i.e. sister group to a given taxon, whether that group is a crown group or not. The term may also mean a group, possibly paraphyletic, defined by primitive traits (i.e. symplesiomorphies). It is generally taken to mean a side branch splitting off earlier on the phylogenetic tree than the group in question.
Palaeontological significance of stem and crown groups
Placing fossils in their right order in a stem group allows the order of these acquisitions to be established, and thus the ecological and functional setting of the evolution of the major features of the group in question. Stem groups thus offer a route to integrate unique palaeontological data into questions of the evolution of living organisms. Furthermore, they show that fossils that were considered to lie in their own separate group because they did not show all the diagnostic features of a living clade, can nevertheless be related to it by lying in its stem group. Such fossils have been of particular importance in considering the origins of the tetrapods, mammals, and animals.
The application of the stem group concept also influenced the interpretation of the organisms of the Burgess shale. Their classification in stem groups to extant phyla, rather than in phyla of their own, is thought by some to make the Cambrian explosion easier to understand without invoking unusual evolutionary mechanisms; however, application of the stem group concept does nothing to ameliorate the difficulties that phylogenetic telescoping poses to evolutionary theorists attempting to understand both macroevolutionary change and the abrupt character of the Cambrian explosion. Overemphasis on the stem group concept threatens to delay or obscure proper recognition of new higher taxa.
Stem groups in systematics
As originally proposed by Karl-Ernst Lauterbach, stem groups should be given the prefix "stem" (i.e. Stem-Aves, Stem-Arthropoda), however the crown group should have no prefix. The latter has not been universally accepted for known groups. A number of paleontologists have opted to apply this approach anyway.
| Biology and health sciences | Phylogenetics and taxonomy | Biology |
1019228 | https://en.wikipedia.org/wiki/Repeating%20crossbow | Repeating crossbow | The repeating crossbow (), also known as the repeater crossbow, and the Zhuge crossbow (, also romanized Chu-ko-nu) due to its association with the Three Kingdoms-era strategist Zhuge Liang (181–234 AD), is a crossbow invented during the Warring States period in China that combined the bow spanning, bolt placing, and shooting actions into one motion.
The earliest archaeological evidence of the repeating crossbow is found in the state of Chu, but it uses a pistol grip that is different from the later and more commonly known Ming dynasty design.
Although the repeating crossbow was in use throughout most of Chinese history until the late Qing dynasty, it was generally regarded as a non-military weapon suited for women, defending households against robbers.
History
According to the Wu-Yue Chunqiu (history of the Wu-Yue War), written in the Eastern Han dynasty, the repeating crossbow was invented during the Warring States Period by a Mr. Qin from the State of Chu. This is corroborated by the earliest archaeological evidence of repeating crossbows, which was excavated from a Chu burial site at Tomb 47 at Qinjiazui, Hubei Province, and has been dated to the 4th century BC, during the Warring States Period (475 - 220 BC). Unlike repeating crossbows of later eras, the ancient double-shot repeating crossbow uses a pistol grip and a rear-pulling mechanism for arming. The Ming repeating crossbow uses an arming mechanism that requires its user to push a rear lever upwards and downwards back and forth. Although handheld repeating crossbows were generally weak and required additional poison, probably aconite, for lethality, much larger mounted versions appeared during the Ming dynasty.
In 180 AD, Yang Xuan used a type of repeating crossbow powered by the movement of wheels:
The invention of the repeating crossbow has often been attributed to Zhuge Liang, but he in fact had nothing to do with it. This misconception is based on a record attributing improvements to the multiple bolt crossbows to him.
During the Ming dynasty, repeating crossbows were used on ships.
Although the repeating crossbow has been used throughout Chinese history and is attested as late as 19th century Qing dynasty in battle against the Japanese, it was generally not regarded as an important military weapon. The Wubei Zhi, written during the 17th century, says that it was favored by people in southeast China but lacked in strength and its bolts tended not to harm anyone. The functions of the repeating crossbow listed in the text are primarily non-military: tiger hunting, defending fortified houses, and usage by timid men and women. According to the Tiangong Kaiwu, also written during the 17th century, the repeating crossbow is only useful against robbers.
Designs
The repeating crossbow combined the actions of spanning the bow, placing the bolt, and shooting into a one-handed movement, thus allowing for a much higher rate of fire than a standard hand-held crossbow. The most common repeating crossbow design originated from the Ming Dynasty and consisted of a top-mounted magazine containing a reservoir of bolts which fed the crossbow using gravity, a rectangular lever attached to both the tiller and the magazine, and a tiller mounting the prods with a stock. By holding the tiller firmly against the hip while pushing and pulling the lever forwards and backwards, the user was able to catch the drawstring on to side notches at the back of the magazine while loading the bolt. A sliding lug nut at the back of the magazine pushed the drawstring out of the notches once the lever is fully pulled backwards; with the tiller pushing the nut up and enabling the drawstring to propel the loaded bolt. The Korean version mounted the magazine at the end of a longer stalk as well as a pivoting recurve bow as a prod; increasing the draw strength, span, range, and performance of the crossbow. Additionally, both the Ming Dynasty in China and the Joseon Dynasty in Korea developed variations that either shot two to three bolts per draw or fired pellets in place of bolts.
An earlier version originated from the State of Chu during the Warring States period used a different design. It consisted of a tiller mounting a fixed double magazine on top as well as a pistol style grip at the bottom beneath the prods mount. Instead of an overhand lever for arming and shooting, it used a sliding lever that had a handle tied to the end with a chord. The lever was pumped forwards and backwards with one hand while the user held the pistol grip firm with the other hand; in a manner similar to drawing a regular bow. Within the crossbow, the lever was embedded with a special metal trigger composed of a latch and sear; the entire trigger being shaped like a crab's claw arm. Upon pushing the lever forward, the trigger was moved forward to catch the drawstring and becomes locked firm by friction and tensional forces from grooves inside the mounting lever and sear. Upon being drawn back, the draw string is spanned while the double magazine fed two bolts onto the firing slots on either side of the trigger once the drawstring is almost fully drawn. At the very end of the pulling action, the sear comes in contact with a round bar that holds the sliding lever in place. The bar pushed the sear forward to release the trigger and enable the drawstring to propel the two loaded bolts. Ultimately, it was superseded by the aforementioned design from the Ming Dynasty due to being overly complex with weaker performance.
Modern design
Another design was more recently manufactured by EK Archery Research, with their Adder Crossbow designed by Jörg Sprave. The crossbow features a magazine to hold five bolts and a self loading arm. The Adder Crossbow, released in 2020, fires at around with a range of
Utility
The basic construction of the repeating crossbow has remained very much unchanged since its invention, making it one of the longest-lived mechanical weapons. The bolts of one magazine are fired and reloaded by simply pushing and pulling the lever back and forth.
The repeating crossbow had an effective range of and a maximum range of . Its comparatively short range limited its usage to primarily defensive positions, where its ability to rapidly fire up to 7–10 bolts in 15–20 seconds was used to prevent assaults on gates and doorways. In comparison, a standard crossbow could only fire about two bolts a minute. The repeating crossbow, with its smaller and lighter ammunition, had neither the power nor the accuracy of a standard crossbow. Thus, it was not very useful against more heavily armoured troops unless poison was smeared on bolts, in which case even a small wound might prove fatal.
| Technology | Archery | null |
1019238 | https://en.wikipedia.org/wiki/Imipramine | Imipramine | Imipramine, sold under the brand name Tofranil, among others, is a tricyclic antidepressant (TCA) mainly used in the treatment of depression. It is also effective in treating anxiety and panic disorder. Imipramine is taken by mouth.
Common side effects of imipramine include dry mouth, drowsiness, dizziness, low blood pressure, rapid heart rate, urinary retention, and electrocardiogram changes. Overdose of the medication can result in death. Imipramine appears to work by increasing levels of serotonin and norepinephrine and by blocking certain serotonin, adrenergic, histamine, and cholinergic receptors.
Imipramine was discovered in 1951 and was introduced for medical use in 1957. It was the first TCA to be marketed. Imipramine and TCAs other than amitriptyline (which, at least in the U.K., is prescribed comparatively as frequently as SSRIs) have decreased in prescription frequency with the rise of SSRIs—which have fewer inherent side effects and are far safer in overdose. Regardless of its caveats, imipramine retains importance in psychopharmacology and pediatrics (e.g., with childhood enuresis).
Medical uses
Imipramine is primarily used for the treatment of depression and certain anxiety disorders, including acute post-traumatic stress reactions. A significant amount of research regarding its efficacy on acute post-traumatic stress in children and adolescents has focused on trauma resulting from burn-injuries.
Although evidence for its efficacy in the treatment of chronic post-traumatic stress disorder appears to be less robust, it remains a viable treatment. Here, it may act fairly similarly to monoamine oxidase inhibitor phenelzine.
Caution is needed in prescribing imipramine (and its commercially-available metabolite, desipramine) in children and youth/adolescents (whether they suffer with, e.g., bed-wetting, recurrent panic attacks, acute trauma or, in the case of desipramine, ADHD), owing to possibility of certain side-effects which may be of particular concern in those under a certain age.
In the treatment of depression, it has demonstrated similar efficacy to the MAOI moclobemide. It has also been used to treat nocturnal enuresis because of its ability to shorten the time of delta wave stage sleep, where wetting occurs. In veterinary medicine, imipramine is used with xylazine to induce pharmacologic ejaculation in stallions. It is also used for separation anxiety in dogs and cats. Blood levels between 150 and 250 ng/mL of imipramine plus its metabolite desipramine generally correspond to antidepressant efficacy.
Contraindications
Combining it with alcohol consumption may cause more drowsiness, necessitating greater caution when drinking. It may be unsafe during pregnancy.
Many MAOIs are known to have serious interactions with imipramine. It is often contraindicated during their use or in the two weeks following their discontinuation. This category includes medications such as isocarboxazid, linezolid, methylene blue, phenelzine, selegiline, moclobemide, procarbazine, rasagiline, safinamide, and tranylcypromine.
Side effects
These side effects can be contributed to the multiple receptors that imipramine targets such as serotonin, norepinephrine, dopamine, acetylcholine, epinephrine, histamine. Those listed in italics below denote common side effects, separated by the organ systems that are affected. Some side effects may be beneficial in some cases, e.g. reduction of hyperactive gag reflex; reduced random or physical strain-linked urinary leakage.
Central nervous system: dizziness, drowsiness, confusion, seizures, headache, anxiety, tremors, stimulation, weakness, insomnia, nightmares, extrapyramidal symptoms in geriatric patients, increased psychiatric symptoms, paresthesia
Cardiovascular: orthostatic hypotension, ECG changes, tachycardia, hypertension, palpitations, dysrhythmias
Eyes, ears, nose and throat: blurred vision, tinnitus, mydriasis
Gastrointestinal: dry mouth, nausea, vomiting, paralytic ileus, increased appetite, cramps, epigastric distress, jaundice, hepatitis, stomatitis, constipation, taste change
Genitourinary: urinary retention
Hematological: agranulocytosis, thrombocytopenia, eosinophilia, leukopenia
Skin: rash, urticaria, diaphoresis, pruritus, photosensitivity
Overdose
Interactions
Like other tricyclic antidepressants, imipramine has many medication interactions. Many MAOIs have serious interactions with this medication. Other categories of medications that may interact with imipramine include blood thinners, antihistamines, muscle relaxants, sleeping pills, thyroid medications, and tranquilizers. Some medications used for various conditions such as high blood pressure, mental illness, nausea, Parkinson's disease, asthma, colds, or allergies.
Certain medications increase the risk of serotonin syndrome, including selective serotonin reuptake inhibitors (SSRIs), St. John's Wort, and drugs such as ecstasy. Other prescription drugs decrease the body's ability to eliminate imipramine. These include barbiturates, some antiarrhythmic medications, some antiepileptic drugs, and certain HIV drugs (protease inhibitors). Others may cause changes in the heart rhythm, such as QT prolongation.
Alcohol and tobacco may interact with imipramine. Tobacco may decrease the medication's effectiveness.
Pharmacology
Pharmacodynamics
Imipramine affects numerous neurotransmitter systems known to be involved in the etiology of depression, anxiety, attention-deficit hyperactivity disorder (ADHD), enuresis and numerous other mental and physical conditions. Imipramine is similar in structure to some muscle relaxants, and has a significant analgesic effect and, thus, is very useful in some pain conditions.
The mechanisms of imipramine's actions include, but are not limited to, effects on:
Serotonin: very strong reuptake inhibition. Imipramine is a tertiary TCA, and is a potent inhibitor of serotonin reuptake, and to a greater extent than secondary amine TCAs such as nortriptyline and despiramine.
Norepinephrine: strong reuptake inhibition. Desipramine has more affinity to norepinephrine transporter than imipramine.
Dopamine: imipramine blocks D2 receptors. Imipramine, and its metabolite desipramine, have no appreciable affinity for the dopamine transporter (Ki = 8,500 and >10,000 nM, respectively).
Acetylcholine: imipramine is, to a certain extent, an antimuscarinic, specifically a relative antagonist of the muscarinic acetylcholine receptors. The attendant side-effects (e.g., blurry vision, dry mouth, constipation), however, are somewhat less common with imipramine than amitriptyline and protriptyline, which tend to cause antimuscarinic side-effects more often. All-in-all, however, it is prescribed with caution to the elderly and with extreme caution to those with psychosis, as the general brain activity enhancement in combination with the "dementing" effects of anticholinergics increases the potential of imipramine to cause hallucinations, confusion and delirium in this population. "Anti-cholinergic" side-effects, including urinary hesitancy/retention, may be treated/reversed with bethanechol and/or other acetylcholine-agonists.
Bethanechol may also be able to alleviate the sexual-dysfunction symptoms which may occur in the context of tricyclic-antidepressant treatment.
Epinephrine: imipramine antagonizes adrenergic receptors, thus sometimes causing orthostatic hypotension.
Sigma receptor: activity on sigma receptors is present, but it is very weak (Ki = 520 nM) and it is about half that of amitriptyline (Ki = 300 nM).
Histamine: imipramine is an antagonist of the histamine H1 receptors.
BDNF: BDNF is implicated in neurogenesis in the hippocampus, and studies suggest that depressed patients have decreased levels of BDNF and reduced hippocampal neurogenesis. It is not clear how neurogenesis restores mood, as ablation of hippocampal neurogenesis in murine models do not show anxiety related or depression related behaviours. Chronic imipramine administration results in increased histone acetylation (which is associated with transcriptional activation and decondensed chromatin) at the hippocampal BDNF promoter, and also reduced expression of hippocampal HDAC5.
Pharmacokinetics
Imipramine has a varied absolute oral bioavailability ranging from 22% to 77%, leading to significant variability in pharmacokinetics. While the drug has rapid and complete absorption after oral administration, much of the drug is affected by first pass metabolism. Food has no effect on absorption, peak drug concentration, or time to peak drug concentration.
Within the body, imipramine is converted into desipramine (desmethylimipramine) as a metabolite.
Chemistry
Imipramine is a tricyclic compound, specifically a dibenzazepine, and possesses three rings fused together with a side chain attached in its chemical structure. Other dibenzazepine TCAs include desipramine (N-desmethylimipramine), clomipramine (3-chloroimipramine), trimipramine (2′-methylimipramine or β-methylimipramine), and lofepramine (N-(4-chlorobenzoylmethyl)desipramine). Imipramine is a tertiary amine TCA, with its side chain-demethylated metabolite desipramine being a secondary amine. Other tertiary amine TCAs include amitriptyline, clomipramine, dosulepin (dothiepin), doxepin, and trimipramine. The chemical name of imipramine is 3-(10,11-dihydro-5H-dibenzo[b,f]azepin-5-yl)-N,N-dimethylpropan-1-amine and its free base form has a chemical formula of C19H24N2 with a molecular weight of 280.407 g/mol. The drug is used commercially mostly as the hydrochloride salt; the embonate (pamoate) salt is used for intramuscular administration and the free base form is not used. The CAS Registry Number of the free base is 50-49-7, of the hydrochloride is 113-52-0, and of the embonate is 10075-24-8.
History
The parent compound of imipramine, 10,11-dihydro-5H-dibenz[b,f]azepine (dibenzazepine), was first synthesized in 1899, but no pharmacological assessment of this compound or any substituted derivatives was undertaken until the late 1940s. Imipramine was first synthesized in 1951, as an antihistamine. The antipsychotic effects of chlorpromazine were discovered in 1952, and imipramine was then developed and studied as an antipsychotic for use in patients with schizophrenia. The medication was tested in several hundred patients with psychosis, but showed little effectiveness. However, imipramine was serendipitously found to possess antidepressant effects in the mid-1950s following a case report of symptom improvement in a woman with severe depression who had been treated with it. This was followed by similar observations in other patients and further clinical research. Subsequently, imipramine was introduced for the treatment of depression in Europe in 1958 and in the United States in 1959. Along with the discovery and introduction of the monoamine oxidase inhibitor iproniazid as an antidepressant around the same time, imipramine resulted in the establishment of monoaminergic drugs as antidepressants.
In the late 1950s, imipramine was the first TCA to be developed (by Ciba). At the first international congress of neuropharmacology in Rome, September 1958 Dr Freyhan from the University of Pennsylvania discussed as one of the first clinicians the effects of imipramine in a group of 46 patients, most of them diagnosed as "depressive psychosis". The patients were selected for this study based on symptoms such as depressive apathy, kinetic retardation and feelings of hopelessness and despair. In 30% of all patients, he reported optimal results, and in around 20%, failure. The side effects noted were atropine-like, and most patients experienced dizziness. Imipramine was first tried for treating psychotic disorders such as schizophrenia, but proved ineffective. As an antidepressant, it did well in clinical studies and it is known to work well in even the most severe cases of depression. It is not surprising, therefore, that imipramine may cause a high rate of manic and hypomanic reactions in hospitalized patients with pre-existing bipolar disorder, with one study showing that up to 25% of such patients maintained on Imipramine switched into mania or hypomania. Such powerful antidepressant properties have made it favorable in the treatment of treatment-resistant depression.
Before the advent of selective serotonin reuptake inhibitors (SSRIs), its sometimes intolerable side-effect profile was considered more tolerable. Therefore, it became extensively used as a standard antidepressant and later served as a prototypical drug for the development of the later-released TCAs. Since SSRIs are superior in terms of inherent side-effect tolerability (although probably inferior in terms of actual efficacy), it has, as of the 1990s, no longer been used as commonly, but is sometimes still prescribed as a second-line treatment for treating major depression. It has also seen limited use in the treatment of migraines, ADHD, and post-concussive syndrome. Imipramine has additional indications for the treatment of panic attacks, chronic pain, and Kleine-Levin syndrome. In pediatric patients, it is relatively frequently used to treat pavor nocturnus and nocturnal enuresis.
Society and culture
Generic names
Imipramine is the English and French generic name of the drug and its , , and , while imipramine hydrochloride is its , , , and . Its generic name in Spanish and Italian and its are imipramina, in German is imipramin, and in Latin is imipraminum. The embonate salt is known as imipramine pamoate.
Brand names
Imipramine is marketed throughout the world mainly under the brand name Tofranil. Imipramine pamoate is marketed under the brand name Tofranil-PM for intramuscular injection.
Availability
Imipramine is available for medical use widely throughout the world, including in the United States, the United Kingdom, elsewhere in Europe, India, Brazil, South Africa, Australia, and New Zealand.
Prescription trends
Between 1998 and 2017, along with amitriptyline, imipramine was the most commonly prescribed first antidepressant for children aged 5-11 years in England.
| Biology and health sciences | Psychiatric drugs | Health |
1019722 | https://en.wikipedia.org/wiki/ISRO | ISRO | The Indian Space Research Organisation (ISRO ) is India's national space agency. It serves as the principal research and development arm of the Department of Space (DoS), overseen by the Prime Minister of India, with the Chairman of ISRO also serving as the chief executive of the DoS. It is primarily responsible for space-based operations, space exploration, international space cooperation and the development of related technologies. The agency maintains a constellation of imaging, communication and remote sensing satellites. It operates the GAGAN and IRNSS satellite navigation systems. It has sent three missions to the Moon and one mission to Mars.
Formerly known as the Indian National Committee for Space Research (INCOSPAR), it was set up in 1962 by then-Prime Minister Jawaharlal Nehru on the recommendation of scientist Vikram Sarabhai. It was renamed as ISRO in 1969 and was subsumed into the Department of Atomic Energy (DAE). The establishment of ISRO institutionalised space research activities in India. In 1972, the Government set up a Space Commission and the DoS, bringing ISRO under its purview. It has since then been managed by the DoS, which also governs various other institutions in the domain of astronomy and space technology.
ISRO built India's first satellite Aryabhata which was launched by the Soviet space agency Interkosmos in 1975. In 1980, it launched the satellite RS-1 onboard the indigenously built launch vehicle SLV-3, making India the seventh country to undertake orbital launches. It has subsequently developed various small-lift and medium-lift launch vehicles, enabling the agency to launch various satellites and deep space missions. It is one of the six government space agencies in the world that possess full launch capabilities with the ability to deploy cryogenic engines, launch extraterrestrial missions and artificial satellites. It is also the only one of four governmental space agencies to have demonstrated unmanned soft landing capabilities.
ISRO's programmes have played a significant role in socio-economic development. It has supported both civilian and military domains in various aspects such as disaster management, telemedicine, navigation and reconnaissance. ISRO's spin-off technologies have also aided in new innovations in engineering and other allied domains.
History
Formative years
Modern space research in India can be traced to the 1920s, when scientist S. K. Mitra conducted a series of experiments sounding the ionosphere through ground-based radio in Kolkata. Later, Indian scientists like C.V. Raman and Meghnad Saha contributed to scientific principles applicable in space sciences. After 1945, important developments were made in coordinated space research in India by two scientists: Vikram Sarabhai, founder of the Physical Research Laboratory at Ahmedabad, and Homi Bhabha, who established the Tata Institute of Fundamental Research in 1945. Initial experiments in space sciences included the study of cosmic radiation, high-altitude and airborne testing, deep underground experimentation at the Kolar mines—one of the deepest mining sites in the world—and studies of the upper atmosphere. These studies were done at research laboratories, universities, and independent locations.
In 1950, the Department of Atomic Energy (DAE) was founded with Bhabha as its secretary. It provided funding for space research throughout India. During this time, tests continued on aspects of meteorology and the Earth's magnetic field, a topic that had been studied in India since the establishment of the Colaba Observatory in 1823. In 1954, the Aryabhatta Research Institute of Observational Sciences (ARIES) was established in the foothills of the Himalayas. The Rangpur Observatory was set up in 1957 at Osmania University, Hyderabad. Space research was further encouraged by the government of India. In 1957, the Soviet Union launched Sputnik 1 and opened up possibilities for the rest of the world to conduct a space launch.
The Indian National Committee for Space Research (INCOSPAR) was set up in 1962 by Prime Minister Jawaharlal Nehru on the suggestion of Dr. Vikram Sarabhai. Initially there was no dedicated ministry for the space programme and all activities of INCOSPAR relating to space technology continued to function within the DAE. IOFS officers were drawn from the Indian Ordnance Factories to harness their knowledge of propellants and advanced light materials used to build rockets. H.G.S. Murthy, an IOFS officer, was appointed the first director of the Thumba Equatorial Rocket Launching Station, where sounding rockets were fired, marking the start of upper atmospheric research in India. An indigenous series of sounding rockets named Rohini was subsequently developed and started undergoing launches from 1967 onwards. Waman Dattatreya Patwardhan, another IOFS officer, developed the propellant for the rockets.
1970's and 1980's
Under the government of Indira Gandhi, INCOSPAR was superseded by ISRO. Later in 1972, a space commission and Department of Space (DoS) were set up to oversee space technology development in India specifically. ISRO was brought under DoS, institutionalising space research in India and forging the Indian space programme into its existing form. India joined the Soviet Interkosmos programme for space cooperation and got its first satellite Aryabhata in orbit through a Soviet rocket.
Efforts to develop an orbital launch vehicle began after mastering sounding rocket technology. The concept was to develop a launcher capable of providing sufficient velocity for a mass of to enter low Earth orbit. It took 7 years for ISRO to develop Satellite Launch Vehicle capable of putting into a orbit. An SLV Launch Pad, ground stations, tracking networks, radars and other communications were set up for a launch campaign. The SLV's first launch in 1979 carried a Rohini technology payload but could not inject the satellite into its desired orbit. It was followed by a successful launch in 1980 carrying a Rohini Series-I satellite, making India the seventh country to reach Earth's orbit after the USSR, the US, France, the UK, China and Japan. RS-1 was the third Indian satellite to reach orbit as Bhaskara had been launched from the USSR in 1979. Efforts to develop a medium-lift launch vehicle capable of putting class spacecrafts into Sun-synchronous orbit had already begun in 1978. They would later lead to the development of the Polar Satellite Launch Vehicle (PSLV). The SLV-3 later had two more launches before discontinuation in 1983. ISRO's Liquid Propulsion Systems Centre (LPSC) was set up in 1985 and started working on a more powerful engine, Vikas, based upon the French Viking. Two years later, facilities to test liquid-fuelled rocket engines were established and development and testing of various rocket engines thrusters began.
At the same time, another solid-fuelled rocket, the Augmented Satellite Launch Vehicle (ASLV), whose design was based upon SLV-3 was being developed, with technologies to launch satellites into geostationary orbit (GTO). The ASLV had limited success and multiple launch failures; it was soon discontinued. Alongside these developments, communication satellite technologies for the Indian National Satellite System and the Indian Remote Sensing Programme for earth observation satellites were developed and launches from overseas were initiated. The number of satellites eventually grew and the systems were established as among the largest satellite constellations in the world, with multi-band communication, radar imaging, optical imaging and meteorological satellites.
1990s
The arrival of the PSLV in 1990s was a major boost for the Indian space programme. With the exception of its first flight in 1994 and two partial failures later, the PSLV had a streak of more than 50 successful flights. The PSLV enabled India to launch all of its low Earth orbit satellites, small payloads to GTO and hundreds of foreign satellites. Along with the PSLV flights, development of a new rocket, a Geosynchronous Satellite Launch Vehicle (GSLV) was going on. India tried to obtain upper-stage cryogenic engines from Russia's Glavkosmos but was blocked by the US from doing so. As a result, KVD-1 engines were imported from Russia under a new agreement which had limited success and a project to develop indigenous cryogenic technology was launched in 1994, taking two decades to reach fulfillment. A new agreement was signed with Russia for seven KVD-1 cryogenic stages and a ground mock-up stage with no technology transfer, instead of five cryogenic stages along with the technology and design in the earlier agreement. These engines were used for the initial flights and were named GSLV Mk.1. ISRO was under US government sanctions between 6 May 1992 to 6 May 1994. After the United States refused to help India with Global Positioning System (GPS) technology during the Kargil war, ISRO was prompted to develop its own satellite navigation system IRNSS (now NaVIC i.e. Navigation with Indian Constellation) which it is now expanding further.
21st century
In 2003, Prime Minister Atal Bihari Vajpayee urged scientists to develop technologies to land humans on the Moon and programmes for lunar, planetary and crewed missions were started. ISRO launched Chandrayaan-1 aboard PSLV in 2008, purportedly the first probe to verify the presence of water on the Moon.
ISRO launched the Mars Orbiter Mission (or Mangalyaan) aboard a PSLV in 2013, which later became the first Asian spacecraft to enter Martian orbit, making India the first country to succeed at this on its first attempt.
Subsequently, the cryogenic upper stage for GSLV rocket became operational, making India the sixth country to have full launch capabilities. A new heavier-lift launcher LVM3 was introduced in 2014 for heavier satellites and future human space missions.
On 23 August 2023, India achieved its first soft landing on an extraterrestrial body and became the first nation to successfully land a spacecraft near the lunar south pole and fourth nation to successfully land a spacecraft on the Moon with ISRO's Chandrayaan-3, the third Moon mission. Indian moon mission, Chandrayaan-3 (lit. "Mooncraft"), saw the successful soft landing of its Vikram lander at 6.04 pm IST (12:34 pm GMT) near the little-explored southern pole of the Moon in a world's first for any space programme.
India then successfully launched its first solar probe, the Aditya-L1, aboard a PSLV on 2 September 2023.
On 30 December 2024, ISRO successfully launched the SpaDeX mission, pioneering spacecraft rendezvous, docking, and undocking using two small satellites. On 16 January 2025, the ISRO Telemetry, Tracking and Command Network's Mission Operations Complex verified that the docking process was successful. India became the 4th country — after USA, Russia and China — to achieve successful Space Docking. ISRO also successfully managed to control two satellites as a single entity after docking.
Agency logo
ISRO has an official logo since 2002. It consists of an orange arrow shooting upwards attached with two blue coloured satellite panels with the name of ISRO written in two sets of text, orange-coloured Devanagari on the left and blue-coloured English in the Prakrta typeface on the right.
Goals and objectives
As the national space agency of India, ISRO's purpose is the pursuit of all space-based applications such as research, reconnaissance, and communications. It undertakes the design and development of space rockets and satellites, and undertakes explores upper atmosphere and deep space exploration missions. ISRO has also incubated technologies in India's private space sector, boosting its growth.
On the topic of the importance of a space programme to India as a developing nation, Vikram Sarabhai as INSCOPAR chair said in 1969:
The former president of India and chairman of DRDO, A. P. J. Abdul Kalam, said:
India's economic progress has made its space programme more visible and active as the country aims for greater self-reliance in space technology. In 2008, India launched as many as 11satellites, including nine from other countries, and went on to become the first nation to launch 10satellites on onerocket. ISRO has put into operation two major satellite systems: the Indian National Satellite System (INSAT) for communication services, and the Indian Remote Sensing Programme (IRS) satellites for management of natural resources.
Organisation structure and facilities
ISRO is managed by the DOS, which itself falls under the authority of the Space Commission and manages the following agencies and institutes:
Indian Space Research Organisation (ISRO)
Antrix Corporation – The marketing arm of ISRO, Bengaluru
Physical Research Laboratory (PRL), Ahmedabad
National Atmospheric Research Laboratory (NARL), Gadanki, Andhra Pradesh
NewSpace India Limited – Commercial wing, Bengaluru
North-Eastern Space Applications Centre (NE-SAC), Umiam
Indian Institute of Space Science and Technology (IIST), Thiruvananthapuram – India's space university
Research facilities
Test facilities
Construction and launch facilities
Tracking and control facilities
Human resource development
Antrix Corporation Limited (Commercial Wing)
Set up as the marketing arm of ISRO, Antrix's job is to promote products, services and technology developed by ISRO.
NewSpace India Limited (Commercial Wing)
Set up for marketing spin-off technologies, tech transfers through industry interface and scale up industry participation in the space programmes.
Space Technology Incubation Centre
ISRO has opened Space Technology Incubation Centres (S-TIC) at premier technical universities in India which will incubate startups to build applications and products in tandem with the industry and for use in future space missions. The S-TIC will bring the industry, academia and ISRO under one umbrella to contribute towards research and development (R&D) initiatives relevant to the Indian Space Programme. S-TICs are at the National Institute of Technology, Agartala serving for east region, National Institute of Technology, Jalandhar for the north region, and the National Institute of Technology, Tiruchirappalli for the south region of India.
Advanced Space Research Group
Similar to NASA's CalTech-operated Jet Propulsion Laboratory, ISRO and the Indian Institute of Space Science and Technology (IIST) implemented a joint working framework in 2021, wherein ISRO will approve all short-, medium- and long-term space research projects of common interest between the two. In return, an Advanced Space Research Group (ASRG) formed at IIST under the guidance of the EOC will have full access to ISRO facilities. This was done with the aim of "transforming" the IIST into a premier space research and engineering institute with the capability of leading future space exploration missions for ISRO.
Directorate of Space Situational Awareness and Management
To reduce dependency on North American Aerospace Defense Command (NORAD) for space situational awareness and protect the civilian and military assets, ISRO is setting up telescopes and radars in four locations to cover each direction. Leh, Mount Abu and Ponmudi were selected to station the telescopes and radars that will cover North, West and South of Indian territory. The last one will be in Northeast India to cover the entire eastern region. Satish Dhawan Space Centre at Sriharikota already supports Multi-Object Tracking Radar (MOTR). All the telescopes and radars will come under Directorate of Space Situational Awareness and Management (DSSAM) in Bengaluru. It will collect tracking data on inactive satellites and will also perform research on active debris removal, space debris modelling and mitigation.
For early warning, ISRO began a ₹400 crore (4 billion; US$53 million) project called Network for Space Object Tracking and Analysis (NETRA). It will help the country track atmospheric entry, intercontinental ballistic missile (ICBM), anti-satellite weapon and other space-based attacks. All the radars and telescopes will be connected through NETRA. The system will support remote and scheduled operations. NETRA will follow the Inter-Agency Space Debris Coordination Committee (IASDCC) and United Nations Office for Outer Space Affairs (UNOSA) guidelines. The objective of NETRA is to track objects at a distance of in GTO.
India signed a memorandum of understanding on the Space Situational Awareness Data Sharing Pact with the US in April 2022. It will enable Department of Space to collaborate with the Combined Space Operation Center (CSpOC) to protect the space-based assets of both nations from natural and man-made threats. On 11 July 2022, ISRO System for Safe and Sustainable Space Operations Management (IS4OM) at Space Situational Awareness Control Centre, in Peenya was inaugurated by Jitender Singh. It will help provide information on on-orbit collision, fragmentation, atmospheric re-entry risk, space-based strategic information, hazardous asteroids, and space weather forecast. IS4OM will safeguard all the operational space assets, identify and monitor other operational spacecraft with close approaches which have overpasses over Indian subcontinent and those which conduct intentional manoeuvres with suspicious motives or seek re-entry within South Asia.
ISRO System for Safe and Sustainable Space Operations Management
On 7 March 2023, ISRO System for Safe and Sustainable Space Operations Management (IS4OM) conducted successful controlled re-entry of decommissioned satellite Megha-Tropiques after firing four on-board 11 Newton thrusters for 20 minutes each. A series of 20 manoeuvres were performed since August 2022 by spending 120 kg fuel. The final telemetry data confirmed disintegtration over Pacific Ocean. It was part of a compliance effort following international guidelines on space debris mitigation.
Speaking at the 42nd annual meeting of the Inter-Agency Space Debris Coordination Committee (IADC) in Bengaluru, S. Somanath stated that the long-term goal is for all Indian space actors—both governmental and non-governmental—to accomplish debris-free space missions by 2030.
Other facilities
Balasore Rocket Launching Station (BRLS) – Balasore
Bhaskaracharya Institute For Space Applications and Geo-Informatics (BISAG), Gandhinagar
Human Space Flight Centre (HSFC), Bengaluru
Indian Regional Navigation Satellite System (IRNSS)
Indian Space Science Data Centre (ISSDC)
Integrated Space Cell
Inter-University Centre for Astronomy and Astrophysics (IUCAA)
ISRO Inertial Systems Unit (IISU) – Thiruvananthapuram
Master Control Facility
National Deep Space Observation Centre (NDSPO)
Regional Remote Sensing Service Centres (RRSSC)
General satellite programmes
Since the launch of Aryabhata in 1975, a number of satellite series and constellations have been deployed by Indian and foreign launchers. At present, ISRO operates one of the largest constellations of active communication and earth imaging satellites for military and civilian uses.
The IRS series
The Indian Remote Sensing satellites (IRS) are India's earth observation satellites. They are the largest collection of remote sensing satellites for civilian use in operation today, providing remote sensing services. All the satellites are placed in polar Sun-synchronous orbit (except GISATs) and provide data in a variety of spatial, spectral and temporal resolutions to enable several programs to be undertaken relevant to national development. The initial versions are composed of the 1 (A, B, C, D) nomenclature while the later versions were divided into sub-classes named based on their functioning and uses including Oceansat, Cartosat, HySIS, EMISAT and ResourceSat etc. Their names were unified under the prefix "EOS" regardless of functioning in 2020. They support a wide range of applications including optical, radar and electronic reconnaissance for Indian agencies, city planning, oceanography and environmental studies.
The INSAT series
The Indian National Satellite System (INSAT) is the country's telecommunication system. It is a series of multipurpose geostationary satellites built and launched by ISRO to satisfy the telecommunications, broadcasting, meteorology and search-and-rescue needs. Since the introduction of the first one in 1983, INSAT has become the largest domestic communication system in the Asia-Pacific Region. It is a joint venture of DOS, the Department of Telecommunications, India Meteorological Department, All India Radio and Doordarshan. The overall coordination and management of INSAT system rests with the Secretary-level INSAT Coordination Committee. The nomenclature of the series was changed to "GSAT" from "INSAT", then further changed to "CMS" from 2020 onwards. These satellites have been used by the Indian Armed Forces as well. GSAT-9 or "SAARC Satellite" provides communication services for India's smaller neighbors.
Gagan Satellite Navigation System
The Ministry of Civil Aviation has decided to implement an indigenous Satellite-Based Regional GPS Augmentation System also known as Space-Based Augmentation System (SBAS) as part of the Satellite-Based Communications, Navigation, Surveillance and Air Traffic Management plan for civil aviation. The Indian SBAS system has been given the acronym GAGAN – GPS Aided GEO Augmented Navigation. A national plan for satellite navigation including implementation of a Technology Demonstration System (TDS) over Indian airspace as a proof of concept has been prepared jointly by Airports Authority of India and ISRO. The TDS was completed during 2007 with the installation of eight Indian Reference Stations at different airports linked to the Master Control Centre located near Bengaluru.
Navigation with Indian Constellation (NavIC)
IRNSS with an operational name NavIC is an independent regional navigation satellite system developed by India. It is designed to provide accurate position information service to users in India as well as the region extending up to from its borders, which is its primary service area. IRNSS provides two types of services, namely, Standard Positioning Service (SPS) and Restricted Service (RS), providing a position accuracy of better than in the primary service area.
Other satellites
Kalpana-1 (MetSat-1) was ISRO's first dedicated meteorological satellite. Indo-French satellite SARAL on 25 February 2013. SARAL (or "Satellite with ARgos and AltiKa") is a cooperative altimetry technology mission, used for monitoring the oceans' surface and sea levels. AltiKa measures ocean surface topography with an accuracy of , compared to on average using altimeters, and with a spatial resolution of .
Launch vehicles
During the 1960s and 1970s, India initiated its own launch vehicles owing to geopolitical and economic considerations. In the 1960s–1970s, the country developed a sounding rocket, and by the 1980s, research had yielded the Satellite Launch Vehicle-3 and the more advanced Augmented Satellite Launch Vehicle (ASLV), complete with operational supporting infrastructure.
Satellite Launch Vehicle
The Satellite Launch Vehicle (known as SLV-3) was the first space rocket to be developed by India. The initial launch in 1979 was a failure followed by a successful launch in 1980 making India the sixth country in world with orbital launch capability. The development of bigger rockets began afterwards.
Augmented Satellite Launch Vehicle
Augmented or Advanced Satellite Launch Vehicle (ASLV) was another small launch vehicle released in 1980s to develop technologies required to place satellites into geostationary orbit. ISRO did not have adequate funds to develop ASLV and PSLV at once. Since ASLV suffered repeated failures, it was dropped in favour of a new project.
Polar Satellite Launch Vehicle
Polar Satellite Launch Vehicle or PSLV is the first medium-lift launch vehicle from India which enabled India to launch all its remote-sensing satellites into Sun-synchronous orbit. PSLV had a failure in its maiden launch in 1993. Besides two other partial failures, PSLV has become the primary workhorse for ISRO with more than 50 launches placing hundreds of Indian and foreign satellites into orbit.
Decade-wise summary of PSLV launches:
Geosynchronous Satellite Launch Vehicle
Geosynchronous Satellite Launch Vehicle was envisaged in 1990s to transfer significant payloads to geostationary orbit. ISRO initially had a great problem realising GSLV as the development of CE-7.5 in India took a decade. The US had blocked India from obtaining cryogenic technology from Russia, leading India to develop its own cryogenic engines.
Decade-wise summary of GSLV Launches:
Launch Vehicle Mark-3
Launch Vehicle Mark-3 (LVM3), previously known as GSLV Mk III, is the heaviest rocket in operational service with ISRO. Equipped with a more powerful cryogenic engine and boosters than GSLV, it has significantly higher payload capacity and allows India to launch all its communication satellites. LVM3 is expected to carry India's first crewed mission to space and will be the testbed for SCE-200 engine which will power India's heavy-lift rockets in the future.
Decade-wise summary of LVM3 launches:
Small Satellite Launch Vehicle
The Small Satellite Launch Vehicle (SSLV) is a small-lift launch vehicle developed by the ISRO with payload capacity to deliver to low Earth orbit () or to Sun-synchronous orbit () for launching small satellites, with the capability to support multiple orbital drop-offs.
Decade-wise summary of SSLV launches:
Rohini Sounding Rockets
Rohini is a series of sounding rockets developed by ISRO for meteorological and atmospheric study. These sounding rockets are capable of carrying payloads of between altitudes of . The ISRO currently uses RH-200, RH-300,Mk-II, RH-560 Mk-II and RH-560 Mk-III rockets, which are launched from the Thumba Equatorial Rocket Launching Station (TERLS) in Thumba and the Satish Dhawan Space Centre in Sriharikota.
Launch facilities
Satish Dhawan Space Centre
Kulasekarapattinam Spaceport
Thumba Equatorial Rocket Launching Station
Human spaceflight programme
The first proposal to send humans into space was discussed by ISRO in 2006, leading to work on the required infrastructure and spacecraft. The trials for crewed space missions began in 2007 with the Space Capsule Recovery Experiment (SRE), launched using the Polar Satellite Launch Vehicle (PSLV) rocket, and safely returned to earth 12 days later.
In 2009, the Indian Space Research Organisation proposed a budget of for its human spaceflight programme. An unmanned demonstration flight was expected after seven years from the final approval and a crewed mission was to be launched after seven years of funding. A crewed mission initially was not a priority and left on the backburner for several years. A space capsule recovery experiment in 2014 and a pad abort test in 2018 were followed by Prime Minister Narendra Modi's announcement in his 2018 Independence Day address that India will send astronauts into space by 2022 on the new Gaganyaan spacecraft. To date, ISRO has developed most of the technologies needed, such as the crew module and crew escape system, space food, and life support systems. The project would cost less than 100 billion (US$1.3 billion) and would include sending two or three Indians to space, at an altitude of , for at least seven days, using a GSLV Mk-III launch vehicle.
Astronaut training and other facilities
The newly established Human Space Flight Centre (HSFC) will coordinate the IHSF campaign. ISRO will set up an astronaut training centre in Bengaluru to prepare personnel for flights in the crewed vehicle. It will use simulation facilities to train the selected astronauts in rescue and recovery operations and survival in microgravity, and will undertake studies of the radiation environment of space. ISRO had to build centrifuges to prepare astronauts for the acceleration phase of the launch. Existing launch facilities at Satish Dhawan Space Centre will have to be upgraded for the Indian human spaceflight campaign. Human Space Flight Centre and Glavcosmos signed an agreement on 1 July 2019 for the selection, support, medical examination and space training of Indian astronauts. An ISRO Technical Liaison Unit (ITLU) was to be set up in Moscow to facilitate the development of some key technologies and establishment of special facilities which are essential to support life in space. Four Indian Air Force personnel finished training at Yuri Gagarin Cosmonaut Training Center in March 2021.
Crewed spacecraft
ISRO is working towards an orbital crewed spacecraft that can operate for seven days in low Earth orbit. The spacecraft, called Gaganyaan, will be the basis of the Indian Human Spaceflight Programme. The spacecraft is being developed to carry up to three people, and a planned upgraded version will be equipped with a rendezvous and docking capability. In its first crewed mission, ISRO's largely autonomous spacecraft will orbit the Earth at altitude for up to seven days with a two-person crew on board. A source in April 2023 suggested that ISRO was aiming for a 2025 launch.
Space station
India plans to build a space station as a follow-up programme to Gaganyaan. ISRO chairman K. Sivan has said that India will not join the International Space Station programme and will instead build a space station on its own. It is expected to be placed in a low Earth orbit at altitude and be capable of harbouring three humans for 1520 days. The rough time-frame is five to seven years after completion of the Gaganyaan project. "Giving out broad contours of the planned space station, Dr. Sivan said it has been envisaged to weigh 20 tonnes and will be placed in an orbit of 400 km above earth where astronauts can stay for 15-20 days. The time frame is 5-7 years after Gaganyaan," he stated.
As per S. Somanath, the Phase1 will be ready by 2028 and the entire space station will be completed by 2035. The space station will be an international platform for collaborative research on future interplanetary missions, microgravity studies, space biology, medicine and research.
Planetary sciences and astronomy
ISRO and Tata Institute of Fundamental Research have operated a balloon launch base at Hyderabad since 1967. Its proximity to the geo-magnetic equator, where both primary and secondary cosmic ray fluxes are low, makes it an ideal location to study diffuse cosmic X-ray background.
ISRO played a role in the discovery of three species of bacteria in the upper stratosphere at an altitude between . The bacteria, highly resistant to ultra-violet radiation, are not found elsewhere on Earth, leading to speculation on whether they are extraterrestrial in origin. They are considered extremophiles, and named as Bacillus isronensis in recognition of ISRO's contribution in the balloon experiments, which led to its discovery, Bacillus aryabhata after India's celebrated ancient astronomer Aryabhata and Janibacter hoylei after the distinguished astrophysicist Fred Hoyle.
Astrosat
Launched in 2015, Astrosat is India's first dedicated multi-wavelength space observatory. Its observation study includes active galactic nuclei, hot white dwarfs, pulsations of pulsars, binary star systems, and supermassive black holes located at the centre of the galaxy.
XPoSat
The X-ray Polarimeter Satellite (XPoSat) is a satellite for studying black holes and polarisation. The spacecraft carries the Polarimeter Instrument in X-rays (POLIX) payload which will study the degree and angle of polarisation of bright astronomical X-ray sources in the energy range 5–30 keV. It launched on 1 January 2024 on a PSLV-DL rocket, and it has an expected operational lifespan of at least five years.
Extraterrestrial exploration
Lunar exploration
Chandryaan () are India's series of lunar exploration spacecraft. The initial mission included an orbiter and controlled impact probe while later missions include landers, rovers and sampling missions.
Chandrayaan-1
Chandrayaan-1 was India's first mission to the Moon. The robotic lunar exploration mission included a lunar orbiter and an impactor called the Moon Impact Probe. ISRO launched it using a modified version of the PSLV on 22 October 2008 from Satish Dhawan Space Centre. It entered lunar orbit on 8 November 2008, carrying high-resolution remote sensing equipment for visible, near infrared, and soft and hard X-ray frequencies. During its 312-day operational period (two years were planned), it surveyed the lunar surface to produce a complete map of its chemical characteristics and three-dimensional topography. The polar regions were of special interest, as they had possible ice deposits. Chandrayaan-1 carried 11 instruments: five Indian and six from foreign institutes and space agencies (including NASA, ESA, the Bulgarian Academy of Sciences, Brown University and other European and North American institutions and companies), which were carried for free. The mission team was awarded the American Institute of Aeronautics and Astronautics SPACE 2009 award, the International Lunar Exploration Working Group's International Co-operation award in 2008, and the National Space Society's 2009 Space Pioneer Award in the science and engineering category.
Chandrayaan-2
Chandrayaan-2, the second mission to the Moon, which included an orbiter, a lander and a rover. It was launched on a Geosynchronous Satellite Launch Vehicle Mark III (GSLV Mk III) on 22 July 2019, consisting of a lunar orbiter, the Vikram lander, and the Pragyan lunar rover, all developed in India. It was the first mission meant to explore the little-explored lunar south pole region. The objective of the Chandrayaan-2 mission was to land a robotic rover to conduct various studies on the lunar surface.
The Vikram lander, carrying the Pragyan rover, was scheduled to land on the near side of the Moon, in the south polar region at a latitude of about 70° S at approximately 1:50 am(IST) on 7 September 2019. However, the lander deviated from its intended trajectory starting from an altitude of , and telemetry was lost seconds before touchdown was expected. A review board concluded that the crash-landing was caused by a software glitch. The lunar orbiter was efficiently positioned in an optimal lunar orbit, extending its expected service time from one year to seven. It was planned that there will be another attempt to soft-land on the Moon in 2023, without an orbiter.
Chandrayaan-3
Chandryaan-3 is India's second attempt to soft-land on the Moon after the partial failure of Chandrayaan-2. The mission only included a lander-rover set and communicated with the orbiter from the previous mission.
On 23 August 2023, ISRO became the first space agency to successfully land a spacecraft near the lunar south pole. ISRO is the fourth space agency ever to land on the Moon.
Mars exploration
Mars Orbiter Mission (MOM) or (Mangalyaan-1)
The Mars Orbiter Mission (MOM), informally known as Mangalyaan (eng: ''MarsCraft'' ) was launched into Earth orbit on 5 November 2013 by the Indian Space Research Organisation (ISRO) and has entered Mars orbit on 24 September 2014. India thus became the first country to have a space probe enter Mars orbit on its first attempt. It was completed at a record low cost of $74 million.
MOM was placed into Mars orbit on 24 September 2014. The spacecraft had a launch mass of , with of five scientific instruments as payload.
The National Space Society awarded the Mars Orbiter Mission team the 2015 Space Pioneer Award in the science and engineering category.
Mars and Moon analogue research station
Researchers from the Birbal Sahni Institute of Palaeosciences (BSIP) and Indian Institute of Science (IISc) have determined that Ladakh is the best site for India's first Mars and Moon analogue research station. The study project is being conducted by BSIP's Binita Phartiyal, IISc's Aloke Kumar who pioneered the idea of building space-bricks from biologically solidified lunar and martian regolith, and Gaganyaan astronaut Shubhanshu Shukla. An analog research station is a location where plans and exercises intended for the Moon and Mars are made. The projected research station would be used for geological and astrobiological research, human studies, crew training, advancing Technology Readiness Levels (TRL), testing space technologies, and engineering integration.
In Ladakh, Aaka Space Studio and ISRO will be leading a 21-day Mars and Moon analog mission. An important step forward in India's efforts to develop human spaceflight and analog research in support of the Gaganyaan program and future missions like Bharatiya Antariksha Station. It will replicate the harsh conditions of extraterrestrial environments. The expedition will test human health and endurance in isolation, acquire biometric data, simulate extraterrestrial landscape, investigate circadian lighting, and test life support technologies. The startup has experimented with technology, human endurance, and habitat design in Rann of Kutch in 2023, simulating lunar conditions.
Solar probes
Aditya-L1
On 2 September 2023, ISRO launched the Aditya-L1 mission to study the solar corona. It is the first Indian space-based solar coronagraph to study the corona in visible and near-infrared bands. The main objective of the mission is to study coronal mass ejections (CMEs), their properties (the structure and evolution of their magnetic fields for example), and consequently constrain parameters that affect space weather. On 6 January 2024, Aditya-L1 spacecraft, India's first solar mission, has successfully entered its final orbit around the first Sun-Earth Lagrangian point (L1), approximately 1.5 million kilometers from Earth.
Future projects
ISRO is developing and operationalising more powerful and less pollutive rocket engines so it can eventually develop much heavier rockets. It also plans km above earth where astronauts can stay for 15–20 days. The time frame is 5–7 years after Gaganyaan, he stated. to develop electric and nuclear propulsion for satellites and spacecraft to reduce their weight and extend their service lives. Long-term plans may include crewed landings on Moon and other planets as well.
Engines and launch vehicles
Semi-cryogenic engine
SCE-200 is a rocket-grade kerosene (dubbed "ISROsene") and liquid oxygen (LOX)-based semi-cryogenic rocket engine inspired by RD-120. The engine will be less polluting and far more powerful. When combined with the LVM3, it will boost its payload capacity; it will be clustered in future to power India's heavy rockets.
Methalox engine
Reusable methane and LOX-based engines are under development. Methane is less pollutive, leaves no residue and hence the engine needs very little refurbishment. The LPSC began cold flow tests of engine prototypes in 2020.
Modular heavy rockets
India's own rockets lack the capacity for launching very heavy satellites to the geostationary orbit beyond 4 ton class, a problem that is planned to be fixed with the introduction of the NGLV.
ISRO is studying heavy (HLV) and super-heavy lift launch vehicles (SHLV). Modular launchers are being designed, with interchangeable parts, to reduce production time. A capacity HLV and an SHLV capable of delivering into orbit have been mentioned in statements and presentations from ISRO officials.
The agency intends to develop a launcher in the 2020s which can carry nearly to geostationary transfer orbit, nearly four times the capacity of the existing LVM3. A rocket family of five medium to heavy-lift class modular rockets described as "Next Generation Launch Vehicle or NGLV" (initially planned as Unified Modular Launch Vehicle or Unified Launch Vehicle) are being planned which will share parts and will replace ISRO's existing PSLV, GSLV and LVM3 rockets completely. The rocket family will be powered by SCE-200 cryogenic engine and will have a capacity of lifting from to to geostationary transfer orbit.
Reusable launch vehicles
There have been two reusable launcher projects ongoing at ISRO. One is the ADMIRE test vehicle, conceived as a VTVL system and another is RLV-TD programme, being run to develop an autonomous spacecraft which will be launched vertically but land like a plane.
To realise a fully re-usable two-stage-to-orbit (TSTO) launch vehicle, a series of technology demonstration missions have been conceived. For this purpose, the winged Reusable Launch Vehicle Technology Demonstrator (RLV-TD) has been configured. The RLV-TD acts as a flying testbed to evaluate various technologies such as hypersonic flight, autonomous landing, powered cruise flight, and hypersonic flight using air-breathing propulsion. First in the series of demonstration trials was the Hypersonic Flight Experiment (HEX). ISRO launched the prototype's test flight, RLV-TD, from the Sriharikota spaceport in February 2016. It weighs around and flew up to a height of . HEX was completed five months later. A scaled-up version of it could serve as fly-back booster stage for the winged TSTO concept. HEX will be followed by a landing experiment (LEX) and return flight experiment (REX).
Spacecraft propulsion and power
Electric thrusters
India has been working on replacing conventional chemical propulsion with Hall-effect and plasma thrusters which would make spacecraft lighter. GSAT-4 was the first Indian spacecraft to carry electric thrusters, but it failed to reach orbit. GSAT-9 launched later in 2017, had xenon-based electric propulsion system for in-orbit functions of the spacecraft. GSAT-20 is expected to be the first fully electric satellite from India.
Alpha source thermoelectric propulsion technology
Radioisotope thermoelectric generator (RTG), also called alpha source thermoelectric technology by ISRO, is a type of atomic battery which uses nuclear decay heat from radioactive material to power the spacecraft. In January 2021, the U R Rao Satellite Centre issued an Expression of Interest (EoI) for design and development of a 100-watt RTG. RTGs ensure much longer spacecraft life and have less mass than solar panels on satellites. Development of RTGs will allow ISRO to undertake long-duration deep space missions to the outer planets.
Radioisotope heater unit
ISRO included two radioisotope heater units developed by the Department of Atomic Energy (DAE) in the propulsion module of Chandrayaan-3 on a trial basis which worked flawlessly.
Nuclear propulsion
ISRO has plans for collaboration with Department of Atomic Energy to power future space missions using nuclear propulsion technology.
Quantum technology
Satellite-based quantum communication
At the Indian Mobile Congress (IMC) 2023, ISRO presented its satellite-based quantum communication technology. It's called quantum key distribution (QKD) technology. According to ISRO, it is creating technologies to thwart quantum computers, which have the ability to readily breach the current generation of encrypted secure communication. A significant milestone for unconditionally secured satellite data communication was reached in September 2023 when ISRO demonstrated free-space quantum communication across a 300-meter distance, including live video conferencing using quantum-key encrypted signals.
Extraterrestrial probes
Lunar exploration
The Lunar Polar Exploration mission (LUPEX) is a planned robotic lunar mission concept by Indian Space Research Organisation (ISRO) and Japan Aerospace Exploration Agency (JAXA) that would send a lunar rover and lander to explore the south pole region of the Moon no earlier than 2026. JAXA is likely to provide the under-development H3 launch vehicle and the rover, while ISRO would be responsible for the lander.
Crewed Lunar Landing
ISRO aims to put an astronaut on the lunar surface by 2040.
Mars exploration
The next Mars mission, Mars Lander Mission or Mangalyaan 2, has been proposed for launch in 2024. The new mission plan includes a rover, helicopter, sky crane and a supersonic parachute.
Venus exploration
ISRO is considering an orbiter mission to Venus called Venus Orbiter Mission, that could launch as early as 2023 to study the planet's atmosphere. Some funds for preliminary studies were allocated in the 2017–18 Indian budget under Space Sciences; solicitations for potential instruments were requested in 2017 and 2018. A mission to Venus is scheduled for 2025 that will include a payload instrument called Venus Infrared Atmospheric Gases Linker (VIRAL) which has been co-developed with the Laboratoire atmosphères, milieux, observations spatiales (LATMOS) under French National Centre for Scientific Research (CNRS) and Roscosmos.
The Venus Orbiter Mission (VOM), which is intended to orbit a spacecraft in the orbit of planet Venus for a better understanding of the Venusian surface and subsurface, atmospheric processes, and influence of Sun on Venusian atmosphere, was approved by the Union Cabinet on September 18, 2024, under the direction of Prime Minister Narendra Modi. Understanding the fundamental processes that have transformed Venus—which is thought to have once been habitable and very comparable to Earth—will be crucial to comprehending the development of Earth and Venus, the sister planets. A total of has been sanctioned for the Venus Orbiter Mission, of which would go toward the spacecraft.
Asteroids and outer solar system
Conceptual studies are underway for spacecraft destined for the asteroids and Jupiter, as well, in the long term. The ideal launch window to send a spacecraft to Jupiter occurs every 33 months. If the mission to Jupiter is launched, a flyby of Venus would be required. Development of RTEG power might allow the agency to further undertake deeper space missions to the other outer planets.
Space telescopes and observatories
AstroSat-2
AstroSat-2 is the successor to the AstroSat mission.
Exoworlds
Exoworlds is a joint proposal by ISRO, IIST and the University of Cambridge for a space telescope dedicated for atmospheric studies of exoplanets, planned for 2025. ExoWorlds is proposed as a dedicated mission for exoplanet spectroscopy in the NUV-VISIBLE-IR ranges. It would be placed in a stable orbit around the earth-sun L2 point.
Indian Spectroscopic and Imaging Space Telescope (INSIST)
The Indian Spectroscopic and Imaging Space Telescope (INSIST) will produce high-resolution deep UV-optical images, and will also have capabilities to carry out low to medium resolution spectroscopy.The INSIST proposal was recommended by ISRO for pre-project phase with seed funding in March 2019.Collaboration with the Canadian Space Agency is also being proposed.
Forthcoming satellites
Geospatial intelligence satellites
A family of 50 artificial intelligence based satellites will be launched by ISRO between 2024 and 2028 to collect geospatial intelligence (GEOINT) in different orbits to track military movements and photograph areas of interest. For the sake of national security, the satellites will monitor the neighboring areas and the international border. It will use thermal, optical, synthetic aperture radar (SAR), among other technologies, for GEOINT application. Each satellite using artificial intelligence will have the ability to communicate and collaborate with the remaining satellites in space at different orbits to monitor the environment for intelligence gathering operations.
Upcoming launch facilities
Kulasekharapatnam Spaceport
Kulasekharapatnam Spaceport is an under-development spaceport in Thoothukudi district of Tamil Nadu. After completion, it would serve as the second launch facility of ISRO. This spaceport will mainly be used by ISRO for launching small payloads.
Applications
Telecommunication
India uses its satellite communication network – one of the largest in the world – for applications such as land management, water resources management, natural disaster forecasting, radio networking, weather forecasting, meteorological imaging and computer communication. Business, administrative services, and schemes such as the National Informatics Centre (NIC) are direct beneficiaries of applied satellite technology.
Military
The Integrated Space Cell, under the Integrated Defence Staff headquarters of the Ministry of Defence, has been set up to utilise more effectively the country's space-based assets for military purposes and to look into threats to these assets. This command will leverage space technology including satellites. Unlike an aerospace command, where the Air Force controls most of its activities, the Integrated Space Cell envisages cooperation and coordination between the three services as well as civilian agencies dealing with space.
With 14 satellites, including GSAT-7A for exclusive military use and the rest as dual-use satellites, India has the fourth largest number of satellites active in the sky which includes satellites for the exclusive use of its air force (IAF) and navy. GSAT-7A, an advanced military communications satellite built exclusively for the Air Force, is similar to the Navy's GSAT-7, and GSAT-7A will enhance the IAF's network-centric warfare capabilities by interlinking different ground radar stations, ground airbases and airborne early warning and control (AWACS) aircraft such as the Beriev A-50 Phalcon and DRDO AEW&CS.
GSAT-7A will also be used by the Army's Aviation Corps for its helicopters and unmanned aerial vehicle (UAV) operations. In 2013, ISRO launched GSAT-7 for the exclusive use of the Navy to monitor the Indian Ocean Region (IOR) with the satellite's 'footprint' and real-time input capabilities to Indian warships, submarines and maritime aircraft. To boost the network-centric operations of the IAF, ISRO launched GSAT-7A in December 2018. The RISAT series of radar-imaging earth observation satellites is also meant for Military use. ISRO launched EMISAT on 1 April 2019. EMISAT is a electronic intelligence (ELINT) satellite. It will improve the situational awareness of the Indian Armed Forces by providing information and the location of hostile radars.
India's satellites and satellite launch vehicles have had military spin-offs. While India's range Prithvi missile is not derived from the Indian space programme, the intermediate range Agni missile is derived from the Indian space programme's SLV-3. In its early years, under Sarabhai and Dhawan, ISRO opposed military applications for its dual-use projects such as the SLV-3. Eventually, the Defence Research and Development Organisation (DRDO)-based missile programme borrowed staff and technology from ISRO. Missile scientist A.P.J. Abdul Kalam (later elected president), who had headed the SLV-3 project at ISRO, took over as missile programme at DRDO. About a dozen scientists accompanied him, helping to design the Agni missile using the SLV-3's solid fuel first stage and a liquid-fuel (Prithvi-missile-derived) second stage. The IRS and INSAT satellites were primarily intended, and used, for civilian-economic applications, but they also offered military spin-offs. In 1996 the Ministry of Defence temporarily blocked the use of IRS-1C by India's environmental and agricultural ministries in order to monitor ballistic missiles near India's borders. In 1997, the Air Force's "Airpower Doctrine" aspired to use space assets for surveillance and battle management.
Academic
Institutions like the Indira Gandhi National Open University and the Indian Institutes of Technology use satellites for educational applications. Between 1975 and 1976, India conducted its largest sociological programme using space technology, reaching 2,400villages through video programming in local languages aimed at educational development via ATS-6 technology developed by NASA. This experiment—named Satellite Instructional Television Experiment (SITE)—conducted large-scale video broadcasts resulting in significant improvement in rural education.
Telemedicine
ISRO has applied its technology for telemedicine, directly connecting patients in rural areas to medical professionals in urban locations via satellite. Since high-quality healthcare is not universally available in some of the remote areas of India, patients in those areas are diagnosed and analysed by doctors in urban centers in real time via video conferencing. The patient is then advised on medicine and treatment, and treated by the staff at one of the 'super-specialty hospitals' per instructions from those doctors. Mobile telemedicine vans are also deployed to visit locations in far-flung areas and provide diagnosis and support to patients.
Biodiversity Information System
ISRO has also helped implement India's Biodiversity Information System, completed in October 2002. Nirupa Sen details the programme: "Based on intensive field sampling and mapping using satellite remote sensing and geospatial modeling tools, maps have been made of vegetation cover on a 1: 250,000 scale. This has been put together in a web-enabled database that links gene-level information of plant species with spatial information in a BIOSPEC database of the ecological hot spot regions, namely northeastern India, Western Ghats, Western Himalayas and Andaman and Nicobar Islands. This has been made possible with collaboration between the Department of Biotechnology and ISRO."
Cartography
The Indian IRS-P5 (CARTOSAT-1) was equipped with high-resolution panchromatic equipment to enable it for cartographic purposes. IRS-P5 (CARTOSAT-1) was followed by a more advanced model named IRS-P6 developed also for agricultural applications. The CARTOSAT-2 project, equipped with single panchromatic camera that supported scene-specific on-spot images, succeeded the CARTOSAT-1 project.
Spin-offs
ISRO's research has been diverted into spin-offs to develop various technologies for other sectors. Examples include bionic limbs for people without limbs, silica aerogel to keep Indian soldiers serving in extremely cold areas warm, distress alert transmitters for accidents, Doppler weather radar and various sensors and machines for inspection work in engineering industries.
International cooperations
ISRO has signed various formal cooperative arrangements in the form of either Agreements or Memoranda of Understanding (MoU) or Framework Agreements with Afghanistan, Algeria, Argentina, Armenia, Australia, Bahrain, Bangladesh, Bolivia, Brazil, Brunei, Bulgaria, Canada, Chile, China, Egypt, Finland, France, Germany, Hungary, Indonesia, Israel, Italy, Japan, Kazakhstan, Kuwait, Maldives, Mauritius, Mexico, Mongolia, Morocco, Myanmar, Norway, Peru, Portugal, South Korea, Russia, São Tomé and Príncipe, Saudi Arabia, Singapore, South Africa, Spain, Oman, Sweden, Syria, Tajikistan, Thailand, Netherlands, Tunisia, Ukraine, United Arab Emirates, United Kingdom, United States, Uzbekistan, Venezuela and Vietnam. Formal cooperative instruments have been signed with international multilateral bodies including European Centre for Medium-Range Weather Forecasts (ECMWF), European Commission, European Organisation for the Exploitation of Meteorological Satellites (EUMETSAT), European Space Agency (ESA) and South Asian Association for Regional Cooperation (SAARC).
Notable collaborative projects
Chandrayaan-1 also carried scientific payloads to the Moon from NASA, ESA, Bulgarian Space Agency, and other institutions/companies in North America and Europe.
For the Gaganyaan mission, ISRO signed a Technical Implementing Plan (TIP) with ESA to provide ground station support.
Indo-French satellite missions
ISRO has two collaborative satellite missions with France's CNES, namely the now retired Megha-Tropiques to study water cycle in the tropical atmosphere and the presently avtive SARAL for altimetry. A third mission consisting of an Earth observation satellite with a thermal infrared imager, TRISHNA (Thermal infraRed Imaging Satellite for High resolution Natural resource Assessment) is being planned by the two countries.
LUPEX
The Lunar Polar Exploration Mission (LUPEX) is a joint Indo-Japanese mission to study the polar surface of the Moon where India is tasked with providing soft landing technologies.
NISAR
NASA-ISRO Synthetic Aperture Radar (NISAR) is a joint Indo-US radar project carrying an L Band and an S Band radar. It will be world's first radar imaging satellite to use dual frequencies.
Some other notable collaborations include:
ISRO operates LUT/MCC under the international COSPAS/SARSAT Programme for Search and Rescue.
India has established a Centre for Space Science and Technology Education in Asia and the Pacific (CSSTE-AP) that is sponsored by the United Nations.
India is a member of the United Nations Committee on the Peaceful Uses of Outer Space, Cospas-Sarsat, International Astronautical Federation, Committee on Space Research (COSPAR), Inter-Agency Space Debris Coordination Committee (IADC), International Space University, and the Committee on Earth Observation Satellite (CEOS).
Contributing to planned BRICS virtual constellation for remote sensing.
Statistics
Last updated: 26 March 2023
Total number of foreign satellites launched by ISRO: 417 (34 countries)
Spacecraft missions: 116
Launch missions: 86
Student satellites: 13
Re-entry missions: 2
Budget for the Department of Space
Corporate affairs
S-band spectrum scam
In India, electromagnetic spectrum, a scarce resource for wireless communication, is auctioned by the Government of India to telecom companies for use. As an example of its value, in 2010, 20 MHz of 3G spectrum was auctioned for . This part of the spectrum is allocated for terrestrial communication (cell phones). However, in January 2005, Antrix Corporation (commercial arm of ISRO) signed an agreement with Devas Multimedia (a private company formed by former ISRO employees and venture capitalists from the US) for lease of S band transponders (amounting to 70 MHz of spectrum) on two ISRO satellites (GSAT 6 and GSAT 6A) for a price of , to be paid over a period of 12 years. The spectrum used in these satellites (2500 MHz and above) is allocated by the International Telecommunication Union specifically for satellite-based communication in India. Hypothetically, if the spectrum allocation is changed for utilisation for terrestrial transmission and if this 70 MHz of spectrum were sold at the 2010 auction price of the 3G spectrum, its value would have been over . This was a hypothetical situation. However, the Comptroller and Auditor-General considered this hypothetical situation and estimated the difference between the prices as a loss to the Indian Government.
There were lapses on implementing official procedures. Antrix/ISRO had allocated the capacity of the above two satellites exclusively to Devas Multimedia, while the rules said it should always be non-exclusive. The Cabinet was misinformed in November 2005 that several service providers were interested in using satellite capacity, while the Devas deal was already signed. Also, the Space Commission was not informed when approving the second satellite (its cost was diluted so that Cabinet approval was not needed). ISRO committed to spending of public money on building, launching, and operating two satellites that were leased out for Devas.
In late 2009, some ISRO insiders exposed information about the Devas-Antrix deal, and the ensuing investigations led to the deal's annulment. G. Madhavan Nair (ISRO Chairperson when the agreement was signed) was barred from holding any post under the Department of Space. Some former scientists were found guilty of "acts of commission" or "acts of omission". Devas and Deutsche Telekom demanded US$2 billion and US$1 billion, respectively, in damages. The Department of Revenue and Ministry of Corporate Affairs began an inquiry into Devas shareholding.
The Central Bureau of Investigation registered a case against the accused in the Antrix-Devas deal under Section 120-B, besides Section 420 of IPC and Section 13(2) read with 13(1)(d) of PC Act, 1988 in March 2015 against the then executive director of Antrix Corporation, two officials of a USA-based company, a Bengaluru-based private multimedia company, and other unknown officials of the Antrix Corporation or the Department of Space.
Devas Multimedia started arbitration proceedings against Antrix in June 2011. In September 2015, the International Court of Arbitration of the International Chamber of Commerce ruled in favour of Devas, and directed Antrix to pay US$672 million (Rs 44.35 billion) in damages to Devas. Antrix opposed the Devas plea for tribunal award in the Delhi High Court.
Heads of ISRO
List of Chairmen (since 1963) of ISRO.
Vikram Sarabhai (1963–1971)
M. G. K. Menon (1972)
Satish Dhawan (1973–1984)
U. R. Rao (1984–1994)
K. Kasturirangan (1994–2003)
G. Madhavan Nair (2003–2009)
K. Radhakrishnan (2009–2014)
Shailesh Nayak (2015)
A. S. Kiran Kumar (2015–2018)
K. Sivan (2018–2022)
S. Somanath (2022–2025)
V. Narayanan (2025–present)
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1020491 | https://en.wikipedia.org/wiki/Dolomite%20%28rock%29 | Dolomite (rock) | Dolomite (also known as dolomite rock, dolostone or dolomitic rock) is a sedimentary carbonate rock that contains a high percentage of the mineral dolomite, CaMg(CO3)2. It occurs widely, often in association with limestone and evaporites, though it is less abundant than limestone and rare in Cenozoic rock beds (beds less than about 66 million years in age). One of the first geologists to distinguish dolomite from limestone was Déodat Gratet de Dolomieu, a French mineralogist and geologist after whom it is named. He recognized and described the distinct characteristics of dolomite in the late 18th century, differentiating it from limestone.
Most dolomite was formed as a magnesium replacement of limestone or of lime mud before lithification. The geological process of conversion of calcite to dolomite is known as dolomitization and any intermediate product is known as dolomitic limestone. The "dolomite problem" refers to the vast worldwide depositions of dolomite in the past geologic record in contrast to the limited amounts of dolomite formed in modern times. Recent research has revealed sulfate-reducing bacteria living in anoxic conditions precipitate dolomite which indicates that some past dolomite deposits may be due to microbial activity.
Dolomite is resistant to erosion and can either contain bedded layers or be unbedded. It is less soluble than limestone in weakly acidic groundwater, but it can still develop solution features (karst) over time. Dolomite rock can act as an oil and natural gas reservoir.
Name
Dolomite takes its name from the 18th-century French mineralogist Déodat Gratet de Dolomieu (1750–1801), who was one of the first to describe the mineral.
The term dolomite refers to both the calcium-magnesium carbonate mineral and to sedimentary rock formed predominantly of this mineral. The term dolostone was introduced in 1948 to avoid confusion between the two. However, the usage of the term dolostone is controversial, because the name dolomite was first applied to the rock during the late 18th century and thus has technical precedence. The use of the term dolostone was not recommended by the Glossary of Geology published by the American Geological Institute.
In old USGS publications, dolomite was referred to as magnesian limestone, a term now reserved for magnesium-deficient dolomites or magnesium-rich limestones.
Description
Dolomite rock is defined as sedimentary carbonate rock composed of more than 50% mineral dolomite. Dolomite is characterized by its nearly ideal 1:1 stoichiometric ratio of magnesium to calcium. It is distinct from high-magnesium limestone in that the magnesium and calcium form ordered layers within the individual dolomite mineral grains, rather than being arranged at random, as they are in high-magnesium calcite grains. In natural dolomite, magnesium is typically between 44 and 50 percent of total magnesium plus calcium, indicating some substitution of calcium into the magnesium layers. A small amount of ferrous iron typically substitutes for magnesium, particularly in more ancient dolomites. Carbonate rock tends to be either almost all calcite or almost all dolomite, with intermediate compositions being quite uncommon.
Dolomite outcrops are recognized in the field by their softness (mineral dolomite has a Mohs hardness of 4 or less, well below common silicate minerals) and because dolomite bubbles feebly when a drop of dilute hydrochloric acid is dropped on it. This distinguishes dolomite from limestone, which is also soft but reacts vigorously with dilute hydrochloric acid. Dolomite usually weathers to a characteristic dull yellow-brown color due to the presence of ferrous iron. This is released and oxidized as the dolomite weathers. Dolomite is usually granular in appearance, with a texture resembling grains of sugar.
Under the microscope, thin sections of dolomite usually show individual grains that are well-shaped rhombs, with considerable pore space. As a result, subsurface dolomite is generally more porous than subsurface limestone and makes up 80% of carbonate rock petroleum reservoirs. This texture contrasts with limestone, which is usually a mixture of grains, micrite (very fine-grained carbonate mud) and sparry cement. The optical properties of calcite and mineral dolomite are difficult to distinguish, but calcite almost never crystallizes as regular rhombs, and calcite is stained by Alizarin Red S while dolomite grains are not. Dolomite rock consisting of well-formed grains with planar surfaces is described as planar or idiotopic dolomite, while dolomite consisting of poorly-formed grains with irregular surfaces is described as nonplanar or xenotopic dolomite. The latter likely forms by recrystallization of existing dolomite at elevated temperature (over ).
The texture of dolomite often shows that it is secondary, formed by replacement of calcium by magnesium in limestone. The preservation of the original limestone texture can range from almost perfectly preserved to completely destroyed. Under a microscope, dolomite rhombs are sometimes seen to replace oolites or skeletal particles of the original limestone. There is sometimes selective replacement of fossils, with the fossil remaining mostly calcite and the surrounding matrix composed of dolomite grains. Sometimes dolomite rhombs are seen cut across the fossil outline. However, some dolomite shows no textural indications that it was formed by replacement of limestone.
Occurrence and origin
Dolomite is widespread in its occurrences, though not as common as limestone. It is typically found in association with limestone or evaporite beds and is often interbedded with limestone. There is no consistent trend in its abundance with age, but most dolomite appears to have formed at high stands of sea level. Little dolomite is found in Cenozoic beds (beds less than 65 million years old), which has been a time of generally low sea levels. Times of high sea level also tend to be times of a greenhouse Earth, and it is possible that greenhouse conditions are the trigger for dolomite formation.
Many dolomites show clear textural indications that they are secondary dolomites, formed by replacement of limestone. However, although much research has gone into understanding this process of dolomitization, the process remains poorly understood. There are also fine-grained dolomites showing no textural indications that they formed by replacement, and it is uncertain whether they formed by replacement of limestone that left no textural traces or are true primary dolomites. This dolomite problem was first recognized over two centuries ago but is still not fully resolved.
The dolomitization reaction
is thermodynamically favorable, with a Gibbs free energy of about -2.2 kcal/mol. In theory, ordinary seawater contains sufficient dissolved magnesium to cause dolomitization. However, because of the very slow rate of diffusion of ions in solid mineral grains at ordinary temperatures, the process can occur only by simultaneous dissolution of calcite and crystallization of dolomite. This in turn requires that large volumes of magnesium-bearing fluids are flushed through the pore space in the dolomitizing limestone. Several processes have been proposed for dolomitization.
The hypersaline model (also known as the evaporative reflux model) is based on the observation that dolomite is very commonly found in association with limestone and evaporites, with the limestone often interbedded with the dolomite. According to this model, dolomitization takes place in a closed basin where seawater is subject to high rates of evaporation. This results in precipitation of gypsum and aragonite, raising the magnesium to calcium ratio of the remaining brine. The brine is also dense, so it sinks into the pore space of any underlying limestone (seepage refluxion), flushing out the existing pore fluid and causing dolomitization. The Permian Basin of North America has been put forward as an example of an environment in which this process took place. A variant of this model has been proposed for sabkha environments in which brine is sucked up into the dolomitizing limestone by evaporation of capillary fluids, a process called evaporative pumping.
Another model is the mixing-zone or Dorag model, in which meteoric water mixes with seawater already present in the pore space, increasing the chemical activity of magnesium relative to calcium and causing dolomitization. The formation of Pleistocene dolomite reefs in Jamaica has been attributed to this process. However, this model has been heavily criticized, with one 2004 review paper describing it bluntly as "a myth". A 2021 paper argued that the mixing zone serves as domain of intense microbial activity which promotes dolomitization.
A third model postulates that normal seawater is the dolomitizing fluid, and the necessary large volumes are flushed through the dolomitizing limestone through tidal pumping. Dolomite formation at Sugarloaf Key, Florida, may be an example of this process. A similar process might occur during rises in sea level, as large volumes of water move through limestone platform rock.
Regardless of the mechanism of dolomitization, the tendency of carbonate rock to be either almost all calcite or almost all dolomite suggests that, once the process is started, it completes rapidly. The process likely occurs at shallow depths of burial, under , where there is an inexhaustible supply of magnesium-rich seawater and the original limestone is more likely to be porous. On the other hand, dolomitization can proceed rapidly at the greater temperatures characterizing deeper burial, if a mechanism exists to flush magnesium-bearing fluids through the beds.
Mineral dolomite has a 12% to 13% smaller volume than calcite per alkali cation. Thus dolomitization likely increases porosity and contributes to the sugary texture of dolomite.
The dolomite problem and primary dolomite
Dolomite is supersaturated in normal seawater by a factor of greater than ten, but dolomite is not seen to precipitate in the oceans. Likewise, geologists have not been successful at precipitating dolomite from seawater at normal temperatures and pressures in laboratory experiments. This is likely due to a very high activation energy for nucleating crystals of dolomite.
The magnesium ion is a relatively small ion, and it acquires a tightly bound hydration shell when dissolved in water. In other words, the magnesium ion is surrounded by a clump of water molecules that are strongly attracted to its positive charge. Calcium is a larger ion and this reduces the strength of binding of its hydration shell, so it is much easier for a calcium ion than a magnesium ion to shed its hydration shell and bind to a growing crystal. It is also more difficult to nucleate a seed crystal of ordered dolomite than disordered high-magnesium calcite. As a result, attempts to precipitate dolomite from seawater precipitate high-magnesium calcite instead. This substance, which has an excess of calcium over magnesium and lacks calcium-magnesium ordering, is sometimes called protodolomite. Raising the temperature makes it easier for magnesium to shed its hydration shell, and dolomite can be precipitated from seawater at temperatures in excess of . Protodolomite also rapidly converts to dolomite at temperatures of or higher. The high temperatures necessary for the formation of dolomite helps explain the rarity of Cenozoic dolomites, since Cenozoic seawater temperatures seldom exceeded 40 °C.
It is possible that microorganisms are capable of precipitating primary dolomite. This was first demonstrated in samples collected at Lagoa Vermelha, Brazil in association with sulfate-reducing bacteria (Desulfovibrio), leading to the hypothesis that sulfate ion inhibits dolomite nucleation. Later laboratory experiments suggest bacteria can precipitate dolomite independently of the sulfate concentration. With time other pathways of interaction between microbial activity and dolomite formation have been added to the discord regarding their role in modulation and generation of polysaccharides, manganese and zinc within the porewater. Meanwhile, a contrary view held by other researchers is that microorganisms precipitate only high-magnesium calcite but leave open the question of whether this can lead to precipitation of dolomite.
Dedolomitization
Dolomitization can sometimes be reversed, and a dolomite bed converted back to limestone. This is indicated by a texture of pseudomorphs of mineral dolomite that have been replaced with calcite. Dedolomitized limestone is typically associated with gypsum or oxidized pyrite, and dedolomitization is thought to occur at very shallow depths through infiltration of surface water with a very high ratio of calcium to magnesium.
Uses
Dolomite is used for many of the same purposes as limestone, including as construction aggregate; in agriculture to neutralize soil acidity and supply calcium and magnesium; as a source of carbon dioxide; as dimension stone; as a filler in fertilizers and other products; as a flux in metallurgy; and in glass manufacturing. It cannot substitute for limestone in chemical processes that require a high-calcium limestone, such as manufacture of sodium carbonate. Dolomite is used for production of magnesium chemicals, such as Epsom salt, and is used as a magnesium supplement. It is also used in the manufacture of refractory materials.
Caves in dolomite rock
As with limestone caves, natural caves and solution tubes typically form in dolomite rock as a result of the dissolution by weak carbonic acid. Caves can also, less commonly, form through dissolution of rock by sulfuric acid. Calcium carbonate speleothems (secondary deposits) in the forms of stalactites, stalagmites, flowstone etc., can also form in caves within dolomite rock. “Dolomite is a common rock type, but a relatively uncommon mineral in speleothems”.
Both the 'Union Internationale de Spéléologie' (UIS) and the American 'National Speleological Society' (NSS), extensively use in their publications, the terms "dolomite" or "dolomite rock" when referring to the natural bedrock containing a high percentage of CaMg(CO3)2 in which natural caves or solution tubes have formed.
Dolomite speleothems
Both calcium and magnesium go into solution when dolomite rock is dissolved. The speleothem precipitation sequence is: calcite, Mg-calcite, aragonite, huntite and hydromagnesite. Hence, the most common speleothem (secondary deposit) in caves within dolomite rock karst, is calcium carbonate in the most stable polymorph form of calcite. Speleothem types known to have a dolomite constituent include: coatings, crusts, moonmilk, flowstone, coralloids, powder, spar and rafts. Although there are reports of dolomite speleothems known to exist in a number of caves around the world, they are usually in relatively small quantities and form in very fine-grained deposits.
| Physical sciences | Sedimentary rocks | Earth science |
1021213 | https://en.wikipedia.org/wiki/Kodiak%20bear | Kodiak bear | The Kodiak bear (Ursus arctos middendorffi), also known as the Kodiak brown bear and sometimes the Alaskan brown bear, inhabits the islands of the Kodiak Archipelago in southwest Alaska. It is one of the largest recognized subspecies or population of the brown bear, and one of the two largest bears alive today, the other being the polar bear. They are also considered by some to be a population of grizzly bears.
Physiologically and physically, the Kodiak bear is very similar to the other brown bear subspecies, such as the mainland grizzly bear (Ursus arctos horribilis) and the extinct California grizzly bear (U. a. californicus), with the main difference being size, as Kodiak bears are on average 1.5 to 2 times larger than their cousins. Despite this large variation in size, the diet and lifestyle of the Kodiak bear do not differ greatly from those of other brown bears.
Kodiak bears have interacted with humans for centuries, especially hunters and other people in the rural coastal regions of the archipelago. The bears are hunted for sport and are encountered by hunters pursuing other species. Less frequently, Kodiak bears are killed by people whose property (such as livestock) or person are threatened. In recent history there has been an increasing focus on conservation and protection of the Kodiak bear population as human activity in its range increases. The IUCN classifies the brown bear (Ursus arctos), of which the Kodiak is a subspecies, as being of "least concern" in terms of endangerment or extinction, though the IUCN does not differentiate between subspecies and thus does not provide a conservation status for the Kodiak population. The Alaska Department of Fish and Game however, along with the United States Fish and Wildlife Service to a lesser extent, closely monitor the size and health of the population and the number of bears hunted in the state.
Description
Taxonomy
Taxonomist C.H. Merriam was the first to recognize the Kodiak bear as a unique subspecies of the brown bear, and he named it "Ursus middendorffi" in honor of the celebrated Baltic naturalist, Dr. A. Th. von Middendorff. Subsequent taxonomic work merged all North American brown bears into a single species (Ursus arctos).
Genetic samples from bears on Kodiak have shown that they are most closely related to brown bears on the Alaska Peninsula and Kamchatka, Russia, and all brown bears roughly north of the US. Kodiak bears have been genetically isolated since at least the last ice age (10,000 to 12,000 years ago) and very little genetic diversity exists within the population. Although the current population is healthy, productive, and has shown no overt adverse signs of inbreeding, it may be more susceptible to new diseases or parasites than other, more diverse brown bear populations.
Color
Hair colors range from pale blonde to orange (typically females or bears from southern parts of the archipelago) to dark brown. Cubs will often retain a white "natal ring" around their neck for the first few years of life. The Kodiak bear's color is similar to that of its close relatives, the mainland American Grizzly bear and Eurasian brown bears.
Size
While there is generally much variation in size between brown bears in different areas, most usually weigh between 115 and 360 kg (254 and 794 lb); the Kodiak bear illustrates island gigantism, commonly reaching sizes of . The size range for females (sows) is from , and for males (boars), it is . Mature males average over the course of the year, and can weigh up to at peak times. Females are typically about 20% smaller and 30% lighter than males, and adult sizes are attained when they are six years old. Bears weigh the least when they emerge from their dens in the spring, and can increase their weight by 20–30% during late summer and autumn. As with other animals, captive Kodiak bears can sometimes weigh considerably more than their wild counterparts.
An average adult male measures in length, and stands tall at the shoulder. The largest recorded wild male weighed , and had a hind foot measurement of . A large male Kodiak bear stands up to tall at the shoulder when it is standing on all four legs. When standing fully upright on its hind legs, a large male could reach a height of . The largest verified size for a captive Kodiak bear was for a specimen that lived at the Dakota Zoo in Bismarck, North Dakota. Nicknamed "Clyde", he weighed when he died in June 1987 at the age of 22. According to zoo director Terry Lincoln, Clyde probably weighed close to a year earlier. He still had a fat layer of when he died..
Kodiak bears are the largest brown bear and are even comparable in size to polar bears. This makes Kodiak bears and polar bears both the two largest members of the bear family and the largest extant terrestrial carnivorans.
The standard method of evaluating the size of bears is by measuring their skulls. Most North American hunting organizations and management agencies use calipers to measure the length of the skull (back of sagittal crest on the back of the skull to the front tooth), and the width (maximum width between the zygomatic arches — "cheek bones"). The total skull size is the sum of these two measurements. The largest bear ever killed in North America was from Kodiak Island, with a total skull size of , and eight of the top 10 brown bears listed in the Boone and Crockett record book are from Kodiak. The average skull size of Kodiak bears that were killed by hunters in the first five years of the 21st century was for boars, and for sows.
Also, an individual named Teddy, which portrayed a killer bear in the movie Grizzly, stood tall on its hind legs and was the largest bear in captivity at the time.
Life history
Distribution and density
This brown bear population only occurs on the islands of the Kodiak Archipelago (Kodiak, Afognak, Shuyak, Raspberry, Uganik, Sitkalidak, and adjacent islands). The Kodiak bear population was estimated to include 3,526 bears in 2005, yielding an estimated archipelago-wide population density of 270 bears per 1000 km2 (700 per 1000 sq. mi). During the past decade, the population has been slowly increasing.
Home range
Bears on Kodiak are naturally active during the day, but when faced with competition for food or space, they adopt a more nocturnal (active at night) lifestyle. This behavior is especially evident in the bears that live near and within Kodiak City. Kodiak bears do not defend territories, but they do have traditional areas that they use each year (home ranges). Because of the rich variety of foods available on Kodiak, the bears on the archipelago have some of the smallest home ranges of any brown bear populations in North America and a great deal of overlap occurs among the ranges of individual bears. Home ranges of adult sows on Kodiak Island average , while boar home ranges average .
Denning
Kodiak bears begin entering their dens in late October. Pregnant sows are usually the first to go to dens; males are the last. Males begin emerging from their dens in early April, while sows with new cubs may stay in dens until late June. Bears living on the north end of Kodiak Island tend to have longer denning periods than bears in the southern areas. Most Kodiak bears dig their dens in hill or mountain sides and they use a wide variety of denning habitats depending on which part of the archipelago they live. Almost a quarter of the adult bears forgo denning, staying somewhat active throughout the winter.
Reproduction and survival
Kodiak bears reach sexual maturity at age five, but most sows are over nine years old when they successfully wean their first litter. The average time between litters is four years. Sows continue to produce cubs throughout their lives, but their productivity diminishes after they are 20 years old. Mating season for Kodiak bears is during May and June. They are serially monogamous (having one partner at a time), staying together from two days to two weeks. As soon as the egg is fertilized and divides a few times, it enters a state of suspended animation until autumn when it finally implants on the uterine wall and begins to grow again. Cubs are born in the den during January or February. Weighing less than at birth with little hair and closed eyes, they suckle for several months, emerging from the den in May or June, weighing . Typical litter sizes on Kodiak are two or three cubs, with a long-term average of 2.4 cubs per litter. However, Kodiak bears have six functional nipples and can litter up to six cubs have been reported. Sows are sometimes seen with five or six cubs in tow, probably due to adopting cubs from other litters. Most cubs stay with their mothers for three years. Almost half of the cubs die before they leave, with cannibalism by adult males being one of the major causes of death.
Kodiak bears that have recently left their mothers, at ages 3–5 years, have high mortality rates with only 56% of males and 89% of females surviving. Most young female bears stay within or near their mother's home range, while most males move farther away. Most adult sows die of natural causes (56%), while most adult male bears are killed by hunters (91%). The oldest known male bear in the wild was 27 years old, and the oldest female was 35.
Habitat
The islands of the Kodiak Archipelago have a subpolar oceanic climate with cool temperatures, overcast skies, fog, windy conditions, and moderate to heavy precipitation throughout most of the year. Although the archipelago only covers about , a rich variety of topography and vegetation ranges from dense forests of Sitka spruce on the northern islands, to steep, glaciated mountains rising to Koniag Peak's along the central spine of Kodiak Island, to rolling hills and flat tundra on the south end of the archipelago. About 14,000 people live on the archipelago, primarily in and around the city of Kodiak and six outlying villages. Roads and other human alterations are generally limited to Afognak Island and the northeastern part of Kodiak Island. About half of the archipelago is included in the Kodiak National Wildlife Refuge.
Feeding habits
Bears live throughout the archipelago, adapting to local resources and retaining relatively small home ranges and comparable densities in most habitats. With such a variety and abundance of food sources, bears are surprisingly intelligent in their eating habits. The first foods bears eat in the spring are emerging vegetation (such as grasses and forbs) and animals that may have died during the winter. This allows the bear to quickly replace the weight that was lost during hibernation. As summer progresses, a wide variety of vegetation supplies nutritional needs until salmon return. Salmon runs extend from May through September on most of the archipelago and bears consume the five species of Pacific salmon that spawn in local streams and lakes. Bears often prioritize the brain, flesh, and eggs of salmon for their high nutritional value. In the late summer and early fall, bears consume several types of berries when they reach their ripest point, and have the highest levels of sugar. As climate change causes elderberries to ripen earlier, berry season is now overlapping with salmon season and some bears are abandoning salmon runs to focus on the berries. Bears also feed on wind-rowed seaweed and invertebrates on some beaches throughout the year. When eating deer, mountain goats, elk, or cattle, internal organs are eaten first for their high-fat content, however even though there is an abundance of the animals found on the archipelago, few Kodiak bears actively prey on them as other methods of finding food are more energy efficient. Another food source available year-round is the garbage made by the human population of Kodiak Island.
Behavior
The Kodiak bear is much like other brown bears in intelligence, although its tendency to feed in large dense groups leads to more complex social behaviors. Kodiak bears are generally solitary in nature; however, when food is concentrated in small areas, such as along salmon spawning streams, grass/sedge flats, berry patches, a dead whale, or even an open garbage dump, they often occur in large groups. Along a few streams on Kodiak, up to 60 bears can be seen simultaneously in a area. To maximize food intake at these ecologically important areas, bears have learned to minimize fighting and fatal interactions by developing a complex communication (both verbal and body posturing) and social structure.
Interactions with people
Usually, Kodiak bears attempt to avoid encounters with people. The most notable exceptions to this behavior pattern occur when bears are surprised, threatened, or attracted by human food, garbage, or hunter-killed game. However, there has been an increase in Kodiak encounters due to increases in the local population as well as increased hunting of Kodiak bears. Bear safety precautions aim at avoiding such situations, understanding bear needs and behavior, and learning how to recognize the warning signs bears give when stressed.
One fatal bear attack on a person on the Kodiak archipelago occurred in 1999. The National Geographic Society filmed a television program about brown bears, which included a segment on two brown bear attacks. Both incidents involved hunters who were hunting by themselves and were returning to game they had killed previously, and left alone in order to continue hunting. One of the attacks was fatal, with the hunter being killed by the bear, and occurred on Uganik Island (November 3, 1999), which is part of the Kodiak archipelago. In the other incident, after being attacked by the bear, the hunter stabbed it with a knife, then recovered his rifle and killed the attacking bear. This occurred on Raspberry Island, home to two full-service wilderness lodges. Prior to that, the last fatality was in 1921. About once every other year, a bear injures a person on Kodiak. In October 2021, a father and son hunting duo survived an attack from a Kodiak bear during an elk registration hunt on Afognak Island.
History and management
Prehistory
Early human occupants of the archipelago when the land was locked into the ice age looked to the sea for their sustenance. At that time, Kodiak Natives (Alutiiqs) occasionally hunted bears, using their meat for food, hides for clothing and bedding, and teeth for adornment. Traditional stories often revolved around the similarity between bears and humans, and the mystical nature of bears because of their proximity to the spirit world.
Commercial harvests
Russian hunters came to the area in the late 18th century to capitalize on the abundant fur resources. Bear hides were considered a "minor fur" and sold for about the same price as river otter pelts. The number of bears harvested increased substantially when sea otter populations declined and after the United States acquired Alaska in 1867, bear harvests on Kodiak increased, peaking at as many as 250 bears per year. Commercial fishing activities increased in the late 1880s and canneries proliferated throughout the archipelago. Bears were viewed as competitors for the salmon resource and were routinely shot when seen on streams or coasts. At the same time, sportsmen and scientists had recognized the Kodiak bear as the largest in the world, and they voiced concerns about overharvesting the population.
Guided hunters and competition for resources
Professional interest in guided Kodiak bear hunts and concern for unregulated resource use in frontier lands such as Alaska prompted the territorial government's newly established Alaska Game Commission to abolish commercial bear hunting (selling the hides) on the archipelago in 1925. The impacts of the new regulations seemed to restore bear populations on the Kodiak Islands. By the 1930s, ranchers in northeast Kodiak reported an increase in bear problems and demanded action. Bears were wrongly seen as a threat to the expanding commercial salmon-fishing industry. To address the dilemma of conserving bears while protecting cattle, salmon, and people, President Franklin D. Roosevelt created the Kodiak National Wildlife Refuge by executive order in 1941. The refuge roughly encompasses the southwestern two-thirds of Kodiak Island, Uganik Island, the Red Peaks area on northwestern Afognak Island, and all of Ban Island.
Alaska achieved statehood in 1959 and assumed responsibility for managing the state's wildlife. The Alaska Board of Game reduced bear-hunting seasons on Afognak and Raspberry Islands and on the Kodiak National Wildlife Refuge, but liberalized bear seasons on non-refuge lands on Kodiak. During the 1960s, state biologists worked with ranchers along the Kodiak road system to examine and reduce the predation problem. Biologists reported that cattle and bears were not compatible on the same ranges and potential solutions included poisons, fences to isolate cattle ranges, and aerial shooting of bears. Again, sport hunters voiced their support for Kodiak bears. Despite public pressure, the state continued actively pursuing and dispatching problem bears until 1970.
Changes in land status
In 1971, the Alaska Native Claims Settlement Act (ANCSA) resolved many long-standing land issues with Aboriginal Alaskans statewide. The impacts were felt strongly on the archipelago as large areas were conveyed to the Native corporations. Federal management of the National Forest lands on Afognak was transferred to Native Corporation ownership with the passage of the Alaska National Interest Lands Conservation Act in 1980 (ANILCA), and the Kodiak National Wildlife Refuge lost control of of prime bear habitat (more than 17% of refuge lands).
In 1975, construction of a logging road began on Afognak Island, and timber harvesting began in 1977. In 1979, work began on an environmental impact statement for the Terror Lake hydroelectric project on Kodiak Island. That project included an earthen dam on Terror Lake with Kodiak National Wildlife Refuge and a tunnel through a mountain ridge to a penstock and powerhouse in the Kizhuyak River drainage. The hydro project was the first significant invasion of inland bear habitat on Kodiak Island. To address the opposition encountered from the public and agencies, a mitigation settlement was negotiated in 1981 which included brown bear research and the establishment of the Kodiak Brown Bear Trust. The hydroelectric project was completed in 1985. Human alteration of bear habitat on Kodiak and Afognak Islands spurred renewed interest and funding for bear research on the archipelago, resulting in a surge of baseline and applied bear research on Kodiak through the 1980s and 1990s.
Bears were not directly harmed by the Exxon Valdez oil spill in 1989, although some were displaced from traditional feeding and traveling areas by cleanup crews. No one was injured by a bear, and no Kodiak bears were killed. To mitigate the adverse impacts of the spill, Exxon reached a settlement with the state and federal governments. Paradoxically, the impacts of the oil spill and the subsequent cleanup and settlement proved to be beneficial to bears on Kodiak. Bear-safety training exposed thousands of workers to factual information about bears, and money from the settlement fund was used for funding land acquisitions. By the close of the 20th century, over 80% of the refuge lands that had been lost as a result of ANCSA and ANILCA were reinstated into the refuge, either through direct purchase or utilizing conservation easements. Lands were also purchased in America, Westtown, and Shuyak Islands and transferred into state ownership. The Kodiak Brown Bear Trust coordinated a coalition of sportsmen and other wildlife conservation groups from around the nation to lobby for the use of settlement funds to acquire Kodiak lands. The groups also directly contributed funding to protect small parcels of important bear habitat around the islands.
Kodiak Archipelago Bear Conservation and Management Plan
In 2001, a citizens advisory committee was established to work closely with the Alaska Department of Fish and Game (ADF&G), with the cooperation of Kodiak NWR, to develop a management plan addressing several problems that affect bears, including hunting, habitat, and viewing. The resulting Kodiak Archipelago Bear Conservation and Management Plan was crafted over several months by representatives from 12 diverse user groups, which, after hearing from a variety of experts from agencies and receiving extensive public input, developed more than 270 recommendations for managing and conserving Kodiak bears. Despite the diversity of viewpoints expressed by members of the group, all of the recommendations were by consensus.
The underlying themes of the recommendations were continued conservation of the bear population at its current level, increased education programs to teach people how to live with bears on Kodiak, and protection of bear habitat with allowances for continued human use of the archipelago. Although the group's role is merely advisory, government management agencies expressed a commitment to implement all of the regulations that were feasible and within their legal jurisdictions.
Genetic diversity and endangerment
The International Union for Conservation of Nature Red List does not list subspecies. The brown bear species, of which the Kodiak subspecies is a member, is listed as Lower Risk or Least Concern. The Kodiak is not listed as an endangered species by the Endangered Species Act of the U.S. Fish and Wildlife Service.
Hunt-management
Kodiak bear research and habitat protection is done cooperatively by the ADF&G and Kodiak National Wildlife Refuge. Bear hunting is managed by the ADF&G, and hunting regulations are established by the Alaska Board of Game. Currently, a finely tuned management system distributes hunters in 32 different areas during two seasons (spring: April 1 – May 15, and fall: October 25 – November 30). Each year, about 4,500 people apply for the 496 permits offered for Kodiak bear hunts (two-thirds to Alaska residents, one-third to non-residents). Nonresidents are required to hire a registered guide who is authorized to hunt in a particular area, and this can cost from $10,000 to $22,000. All hunters must come into the Alaska Department of Fish and Game office in Kodiak before going into the field for a brief orientation and must check out before they leave the island. Every bear that is legally killed on the archipelago must be inspected by an ADF&G wildlife biologist before it can be taken from the islands. Pelts receive a stamp from an ADF&G officer if the hunter and guide provide proper documentation to prove licensing. Pelts cannot be transported or legally preserved or sold without the official stamp. Hunting laws are strictly enforced by the ADF&G officers who often have the full support of the local community. Illegal hunting and fishing is frowned upon by the community which maintains a healthy respect for the island's environmental laws, as well. Stiff penalties accompany illegal hunting and fishing. The island's remote location makes trafficking in illegal pelts difficult for would-be poachers.
Since statehood, the reported number of Kodiak bears killed by hunters has ranged from 77 (1968–1969) to 206 (1965–1966). From 2000 to 2006, an average of 173 Kodiak bears were killed by hunters each year (118 during the fall season and 55 in the spring season). Over 75% of those were males. An additional nine bears were reported killed annually in defense of life or property during the same time. The number of large, trophy-sized bears (total skull size at least ) killed by hunters in recent years has been increasing. In the 1970s, only 2.5% of the bears killed on Kodiak were trophy-sized; in the 1990s and 2000s, the proportion increased to almost 9%.
Bear-viewing
In the past 20 years, bear viewing has become increasingly popular on Kodiak and other parts of Alaska. The most accessible bear-viewing location on Kodiak, Frazer River, had over 1,100 people come in 2007. Visitor numbers have been increasing at about 10% annually and development of additional bear viewing areas on Kodiak is planned. Also, other bear viewing opportunities exist through air-taxi, charter boat, remote lodge, and trekking operations on the archipelago. Although bear-viewing is often considered a "non-consumptive" use, it can have serious impacts on bear populations if it is not conducted properly. Most viewing occurs at places where bears congregate because of feeding opportunities that are critical to their survival. If some bears avoid these areas because people are there, those bears may not get the fat and protein they need to make it through the upcoming winter. Consequently, unmanaged bear viewing could affect several bears, especially productive sows with cubs. Often, bear-viewing and bear-hunting are considered incompatible. Even if the bear population is healthy and bear hunting is sustainable, ethical questions arise especially if hunting occurs near viewing areas and either during or soon after the viewing season. Many feel that it is not fair to encourage bears to be close to people during the summer, only to allow them to be shot in the fall. The Kodiak bear plan recognized bear hunting as a legitimate, traditional, and biologically justifiable activity. It recommended that agencies find ways to make bear hunting and bear viewing compatible on the archipelago.
Cultural significance
The bear is important to the Alutiiq people. Its Alutiiq name is Taquka’aq (Bear), with the pronunciation varying between Northern and Southern dialects.
| Biology and health sciences | Bears | Animals |
1021764 | https://en.wikipedia.org/wiki/Grasshopper | Grasshopper | Grasshoppers are a group of insects belonging to the suborder Caelifera. They are amongst what are possibly the most ancient living groups of chewing herbivorous insects, dating back to the early Triassic around 250 million years ago.
Grasshoppers are typically ground-dwelling insects with powerful hind legs which allow them to escape from threats by leaping vigorously. Their front legs are shorter and used for grasping food. As hemimetabolous insects, they do not undergo complete metamorphosis; they hatch from an egg into a nymph or "hopper" which undergoes five moults, becoming more similar to the adult insect at each developmental stage. The grasshopper hears through the tympanal organ which can be found in the first segment of the abdomen attached to the thorax; while its sense of vision is in the compound eyes, a change in light intensity is perceived in the simple eyes (ocelli). At high population densities and under certain environmental conditions, some grasshopper species can change colour and behavior and form swarms. Under these circumstances, they are known as locusts.
Grasshoppers are plant-eaters, with a few species at times becoming serious pests of cereals, vegetables and pasture, especially when they swarm in the millions as locusts and destroy crops over wide areas. They protect themselves from predators by camouflage; when detected, many species attempt to startle the predator with a brilliantly coloured wing flash while jumping and (if adult) launching themselves into the air, usually flying for only a short distance. Other species such as the rainbow grasshopper have warning coloration which deters predators. Grasshoppers are affected by parasites and various diseases, and many predatory creatures feed on both nymphs and adults. The eggs are subject to attack by parasitoids and predators. Grasshoppers are diurnal insects, meaning they are most active during the day time.
Grasshoppers have had a long relationship with humans. Swarms of locusts can have devastating effects and cause famine, having done so since Biblical times. Even in smaller numbers, the insects can be serious pests. They are used as food in countries such as Mexico and Indonesia. They feature in art, symbolism and literature. The study of grasshopper species is called acridology.
Phylogeny
Grasshoppers belong to the suborder Caelifera. Although "grasshopper" has been used as a common name for the suborder in general, modern sources restrict it to the more "evolved" families. They may be placed in the infraorder Acrididea and have been referred to as "short-horned grasshoppers" in older texts to distinguish them from the also-obsolete term "long-horned grasshoppers" (now bush-crickets or katydids) with their much longer antennae. The phylogeny of the Caelifera, based on mitochondrial ribosomal RNA of thirty-two taxa in six out of seven superfamilies, is shown as a cladogram. The Ensifera (crickets, etc.), Caelifera and all the superfamilies of grasshoppers except "Pamphagoidea" appear to be monophyletic.
In evolutionary terms, the split between the Caelifera and the Ensifera is no more recent than the Permo-Triassic boundary; the earliest insects that are certainly Caeliferans are in the extinct families Locustopseidae and Locustavidae from the early Triassic, roughly 250 million years ago. The group diversified during the Triassic and have remained important plant-eaters from that time to now. The first modern families such as the Eumastacidae, Tetrigidae and Tridactylidae appeared in the Cretaceous, though some insects that might belong to the last two of these groups are found in the early Jurassic. Morphological classification is difficult because many taxa have converged towards a common habitat type; recent taxonomists have concentrated on the internal genitalia, especially those of the male. This information is not available from fossil specimens, and the palaeontological taxonomy is founded principally on the venation of the hindwings.
The Caelifera includes some 2,400 valid genera and about 11,000 known species. Many undescribed species probably exist, especially in tropical wet forests. The Caelifera have a predominantly tropical distribution with fewer species known from temperate zones, but most of the superfamilies have representatives worldwide. They are almost exclusively herbivorous and are probably the oldest living group of chewing herbivorous insects.
The most diverse superfamily is the Acridoidea, with around 8,000 species. The two main families in this are the Acrididae (grasshoppers and locusts) with a worldwide distribution, and the Romaleidae (lubber grasshoppers), found chiefly in the New World. The Ommexechidae and Tristiridae are South American, and the Lentulidae, Lithidiidae and Pamphagidae are mainly African. The Pauliniids are nocturnal and can swim or skate on water, and the Lentulids are wingless. Pneumoridae are native to Africa, particularly southern Africa, and are distinguished by the inflated abdomens of the males.
Characteristics
Grasshoppers have the typical insect body plan of head, thorax, and abdomen. The head is held vertically at an angle to the body, with the mouth at the bottom. The head bears a large pair of compound eyes which give all-round vision, three simple eyes which can detect light and dark, and a pair of thread-like antennae that are sensitive to touch and smell. The downward-directed mouthparts are modified for chewing and there are two sensory palps in front of the jaws.
The thorax and abdomen are segmented and have a rigid cuticle made up of overlapping plates composed of chitin. The three fused thoracic segments bear three pairs of legs and two pairs of wings. The forewings, known as tegmina, are narrow and leathery while the hindwings are large and membranous, the veins providing strength. The legs are terminated by claws for gripping. The hind leg is particularly powerful. The legs of these species are so powerful that they can jump quite a long distance. they also use this to flee from danger. The femur is robust and has several ridges where different surfaces join and the inner ridges bear stridulatory pegs in some species. The posterior edge of the tibia bears a double row of spines and there are a pair of articulated spurs near its lower end. The interior of the thorax houses the muscles that control the wings and legs.
The abdomen has eleven segments, the first of which is fused to the thorax and contains the tympanal organ and hearing system. Segments two to eight are ring-shaped and joined by flexible membranes. Segments nine to eleven are reduced in size; segment nine bears a pair of cerci and segments ten and eleven have the reproductive organs. Female grasshoppers are normally larger than males, with short ovipositors. The name of the suborder "Caelifera" comes from the Latin and means chisel-bearing, referring to the shape of the ovipositor.
The grasshopper's auditory organs are located on its abdomen, rather than on its head. These organs consist of a pair of membranes, each positioned on either side of the first abdominal segment and tucked under the wings. Known as tympanal organs, these simple eardrums vibrate in response to sound waves, enabling the grasshopper to hear the songs of other grasshoppers.
Those species that make easily heard noises usually do so by rubbing a row of pegs on the hind legs against the edges of the forewings (stridulation). These sounds are produced mainly by males to attract females, though in some species the females also stridulate.
Grasshoppers may be confused with crickets, but they differ in many aspects; these include the number of segments in their antennae and the structure of the ovipositor, as well as the location of the tympanal organ and the methods by which sound is produced. Ensiferans have antennae that can be much longer than the body and have at least 20–24 segments, while caeliferans have fewer segments in their shorter, stouter antennae.
Biology
Diet and digestion
Most grasshoppers are polyphagous, eating vegetation from multiple plant sources, but some are omnivorous and also eat animal tissue and animal faeces. In general their preference is for grasses, including many cereals grown as crops. The digestive system is typical of insects, with Malpighian tubules discharging into the midgut. Carbohydrates are digested mainly in the crop, while proteins are digested in the ceca of the midgut. Saliva is abundant but largely free of enzymes, helping to move food and Malpighian secretions along the gut. Some grasshoppers possess cellulase, which by softening plant cell walls makes plant cell contents accessible to other digestive enzymes. Grasshoppers can also be cannibalistic when swarming.
Sensory organs
Grasshoppers have a typical insect nervous system, and have an extensive set of external sense organs. On the side of the head are a pair of large compound eyes which give a broad field of vision and can detect movement, shape, colour and distance. There are also three simple eyes (ocelli) on the forehead which can detect light intensity, a pair of antennae containing olfactory (smell) and touch receptors, and mouthparts containing gustatory (taste) receptors. At the front end of the abdomen there is a pair of tympanal organs for sound reception. There are numerous fine hairs (setae) covering the whole body that act as mechanoreceptors (touch and wind sensors), and these are most dense on the antennae, the palps (part of the mouth), and on the cerci at the tip of the abdomen. There are special receptors (campaniform sensillae) embedded in the cuticle of the legs that sense pressure and cuticle distortion. There are internal "chordotonal" sense organs specialized to detect position and movement about the joints of the exoskeleton. The receptors convey information to the central nervous system through sensory neurons, and most of these have their cell bodies located in the periphery near the receptor site itself.
Circulation and respiration
Like other insects, grasshoppers have an open circulatory system and their body cavities are filled with haemolymph. A heart-like structure in the upper part of the abdomen pumps the fluid to the head from where it percolates past the tissues and organs on its way back to the abdomen. This system circulates nutrients throughout the body and carries metabolic wastes to be excreted into the gut. Other functions of the haemolymph include wound healing, heat transfer and the provision of hydrostatic pressure, but the circulatory system is not involved in gaseous exchange. Respiration is performed using tracheae, air-filled tubes, which open at the surfaces of the thorax and abdomen through pairs of valved spiracles. Larger insects may need to actively ventilate their bodies by opening some spiracles while others remain closed, using abdominal muscles to expand and contract the body and pump air through the system.
Jumping
Grasshoppers jump by extending their large back legs and pushing against the substrate (the ground, a twig, a blade of grass or whatever else they are standing on); the reaction force propels them into the air. A large grasshopper, such as a locust, can jump about a metre (20 body lengths) without using its wings; the acceleration peaks at about 20 g.
They jump for several reasons; to escape from a predator, to launch themselves into flight, or simply to move from place to place. For the escape jump in particular there is strong selective pressure to maximize take-off velocity, since this determines the range. This means that the legs must thrust against the ground with both high force and a high velocity of movement. A fundamental property of muscle is that it cannot contract with high force and high velocity at the same time. Grasshoppers overcome this by using a catapult mechanism to amplify the mechanical power produced by their muscles.
The jump is a three-stage process. First, the grasshopper fully flexes the lower part of the leg (tibia) against the upper part (femur) by activating the flexor tibiae muscle (the back legs of the grasshopper in the top photograph are in this preparatory position). Second, there is a period of co-contraction in which force builds up in the large, pennate extensor tibiae muscle, but the tibia is kept flexed by the simultaneous contraction of the flexor tibiae muscle. The extensor muscle is much stronger than the flexor muscle, but the latter is aided by specialisations in the joint that give it a large effective mechanical advantage over the former when the tibia is fully flexed. Co-contraction can last for up to half a second, and during this period the extensor muscle shortens and stores elastic strain energy by distorting stiff cuticular structures in the leg. The extensor muscle contraction is quite slow (almost isometric), which allows it to develop high force (up to 14 N in the desert locust), but because it is slow only low power is needed. The third stage of the jump is the trigger relaxation of the flexor muscle, which releases the tibia from the flexed position. The subsequent rapid tibial extension is driven mainly by the relaxation of the elastic structures, rather than by further shortening of the extensor muscle. In this way the stiff cuticle acts like the elastic of a catapult, or the bow of a bow-and-arrow. Energy is put into the store at low power by slow but strong muscle contraction, and retrieved from the store at high power by rapid relaxation of the mechanical elastic structures.
Stridulation
Male grasshoppers spend much of the day stridulating, singing more actively under optimal conditions and being more subdued when conditions are adverse; females also stridulate, but their efforts are insignificant when compared to the males. Late-stage male nymphs can sometimes be seen making stridulatory movements, although they lack the equipment to make sounds, demonstrating the importance of this behavioural trait. The songs are a means of communication; the male stridulation seems to express reproductive maturity, the desire for social cohesion and individual well-being. Social cohesion becomes necessary among grasshoppers because of their ability to jump or fly large distances, and the song can serve to limit dispersal and guide others to favourable habitat. The generalised song can vary in phraseology and intensity, and is modified in the presence of a rival male, and changes again to a courtship song when a female is nearby. In male grasshoppers of the family Pneumoridae, the enlarged abdomen amplifies stridulation.
Life cycle
In most grasshopper species, conflicts between males over females rarely escalate beyond ritualistic displays. Some exceptions include the chameleon grasshopper (Kosciuscola tristis), where males may fight on top of ovipositing females; engaging in leg grappling, biting, kicking and mounting.
The newly emerged female grasshopper has a preoviposition period of a week or two while she increases in weight and her eggs mature. After mating, the female of most species digs a hole with her ovipositor and lays a batch of eggs in a pod in the ground near food plants, generally in the summer. After laying the eggs, she covers the hole with soil and litter. Some, like the semi-aquatic Cornops aquaticum, deposit the pod directly into plant tissue. The eggs in the pod are glued together with a froth in some species. After a few weeks of development, the eggs of most species in temperate climates go into diapause, and pass the winter in this state. Diapause is broken by a sufficiently low ground temperature, with development resuming as soon as the ground warms above a certain threshold temperature. The embryos in a pod generally all hatch out within a few minutes of each other. They soon shed their membranes and their exoskeletons harden. These first instar nymphs can then jump away from predators.
Grasshoppers undergo incomplete metamorphosis: they repeatedly moult, each instar becoming larger and more like an adult, with the wing-buds increasing in size at each stage. The number of instars varies between species but is often six. After the final moult, the wings are inflated and become fully functional. The migratory grasshopper, Melanoplus sanguinipes, spends about 25 to 30 days as a nymph, depending on sex and temperature, and lives for about 51 days as an adult.
Swarming
Locusts are the swarming phase of certain species of short-horned grasshoppers in the family Acrididae. Swarming behaviour is a response to overcrowding. Increased tactile stimulation of the hind legs causes an increase in levels of serotonin. This causes the grasshopper to change colour, feed more and breed faster. The transformation of a solitary individual into a swarming one is induced by several contacts per minute over a short period.
Following this transformation, under suitable conditions dense nomadic bands of flightless nymphs known as "hoppers" can occur, producing pheromones which attract the insects to each other. With several generations in a year, the locust population can build up from localised groups into vast accumulations of flying insects known as plagues, devouring all the vegetation they encounter. The largest recorded locust swarm was one formed by the now-extinct Rocky Mountain locust in 1875; the swarm was long and wide, and one estimate puts the number of locusts involved at 3.5 trillion. An adult desert locust can eat about of plant material each day, so the billions of insects in a large swarm can be very destructive, stripping all the foliage from plants in an affected area and consuming stems, flowers, fruits, seeds and bark.
Predators, parasites, and pathogens
Grasshoppers have a wide range of predators at different stages of their lives; eggs are eaten by bee-flies, ground beetles and blister beetles; hoppers and adults are taken by other insects such as ants and , robber flies and sphecid wasps, by spiders, and by many birds and small mammals including dogs and cats.
The eggs and nymphs are under attack by parasitoids including blow flies, flesh flies, and tachinid flies. External parasites of adults and nymphs include mites. Female grasshoppers parasitised by mites produce fewer eggs and thus have fewer offspring than unaffected individuals.
The grasshopper nematode (Mermis nigrescens) is a long slender worm that infects grasshoppers, living in the insects' hemocoel. Adult worms lay eggs on plants and the host becomes infected when the foliage is eaten. Spinochordodes tellinii and Paragordius tricuspidatus are parasitic worms that infect grasshoppers and alter the behaviour of their hosts. When the worms are sufficiently developed, the grasshopper is persuaded to leap into a nearby body of water where it drowns, thus enabling the parasite to continue with the next stage of its life cycle, which takes place in water.
Grasshoppers are affected by diseases caused by bacteria, viruses, fungi and protozoa. The bacteria Serratia marcescens and Pseudomonas aeruginosa have both been implicated in causing disease in grasshoppers, as has the entomopathogenic fungus Beauveria bassiana. This widespread fungus has been used to control various pest insects around the world, but although it infects grasshoppers, the infection is not usually lethal because basking in the sun has the result of raising the insects' temperature above a threshold tolerated by the fungus. The fungal pathogen Entomophaga grylli is able to influence the behaviour of its grasshopper host, causing it to climb to the top of a plant and cling to the stem as it dies. This ensures wide dispersal of the fungal spores liberated from the corpse.
The fungal pathogen Metarhizium acridum is found in Africa, Australia and Brazil where it has caused epizootics in grasshoppers. It is being investigated for possible use as a microbial insecticide for locust control. The microsporidian fungus Nosema locustae, once considered to be a protozoan, can be lethal to grasshoppers. It has to be consumed by mouth and is the basis for a bait-based commercial microbial pesticide. Various other microsporidians and protozoans are found in the gut.
Anti-predator defences
Grasshoppers exemplify a range of anti-predator adaptations, enabling them to avoid detection, to escape if detected, and in some cases to avoid being eaten if captured. Grasshoppers are often camouflaged to avoid detection by predators that hunt by sight; some species can change their coloration to suit their surroundings.
Several species such as the hooded leaf grasshopper Phyllochoreia ramakrishnai (Eumastacoidea) are detailed mimics of leaves. Stick grasshoppers (Proscopiidae) mimic wooden sticks in form and coloration. Grasshoppers often have deimatic patterns on their wings, giving a sudden flash of bright colours that may startle predators long enough to give time to escape in a combination of jump and flight.
Some species are genuinely aposematic, having both bright warning coloration and sufficient toxicity to dissuade predators. Dictyophorus productus (Pyrgomorphidae) is a "heavy, bloated, sluggish insect" that makes no attempt to hide; it has a bright red abdomen. A Cercopithecus monkey that ate other grasshoppers refused to eat the species. Another species, the rainbow or painted grasshopper of Arizona, Dactylotum bicolor (Acridoidea), has been shown by experiment with a natural predator, the little striped whiptail lizard, to be aposematic.
Relationship with humans
In art and media
Grasshoppers are occasionally depicted in artworks, such as the Dutch Golden Age painter Balthasar van der Ast's still life oil painting, Flowers in a Vase with Shells and Insects, c. 1630, now in the National Gallery, London, though the insect may be a bush-cricket.
Another orthopteran is found in Rachel Ruysch's still life Flowers in a Vase, c. 1685. The seemingly static scene is animated by a "grasshopper on the table that looks about ready to spring", according to the gallery curator Betsy Wieseman, with other invertebrates including a spider, an ant, and two caterpillars.
Grasshoppers are also featured in cinema. The 1957 film Beginning of the End portrayed giant grasshoppers attacking Chicago. In the 1998 Disney/Pixar animated film A Bug's Life, the antagonists are a gang of grasshoppers, with their leader Hopper serving as the main villain.
The protagonists of the 1971 tokusatsu series Kamen Rider primarily carry a grasshopper motif (for example Kamen Rider Black's Batta Man form), which continues to serve as the baseline visual template for most entries in the media franchise it has given birth to since.
Symbolism
Grasshoppers are sometimes used as symbols. During the Greek Archaic Era, the grasshopper was the symbol of the polis of Athens, possibly because they were among the most common insects on the dry plains of Attica. Native Athenians for a while wore golden grasshopper brooches to symbolise that they were of pure Athenian lineage with no foreign ancestors. In addition, Peisistratus hung the figure of a kind of grasshopper before the Acropolis of Athens as apotropaic magic.
Another symbolic use of the grasshopper is Sir Thomas Gresham's gilded grasshopper in Lombard Street, London, dating from 1563; the building was for a while the headquarters of the Guardian Royal Exchange, but the company declined to use the symbol for fear of confusion with the locust.
Grasshoppers appearing in dreams have been interpreted as symbols of "Freedom, independence, spiritual enlightenment, inability to settle down or commit to decision". Locusts are taken literally to mean devastation of crops in the case of farmers; figuratively as "wicked men and women" for non-farmers; and "Extravagance, misfortune, & ephemeral happiness" by "gypsies".
As food
In some countries, grasshoppers are used as food. In southern Mexico, grasshoppers, known as chapulines, are eaten in a variety of dishes, such as in tortillas with chilli sauce. Grasshoppers are served on skewers in some Chinese food markets, like the Donghuamen Night Market. Fried grasshoppers (walang goreng) are eaten in the Gunung Kidul Regency, Yogyakarta, Java in Indonesia. Grasshoppers are a beloved delicacy in Uganda; they are usually eaten fried (most commonly in November and May after the rains). In America, the Ohlone burned grassland to herd grasshoppers into pits where they could be collected as food.
It is recorded in the Bible that John the Baptist ate locusts and wild honey (Greek: ἀκρίδες καὶ μέλι ἄγριον, akrídes kaì méli ágrion) while living in the wilderness. However, because of a tradition of depicting him as an ascetic, attempts have been made to explain that the locusts were in fact a suitably ascetic vegetarian food such as carob beans, notwithstanding the fact that the word ἀκρίδες means plainly grasshoppers.
In recent years, with the search for alternative healthy and sustainable protein sources, grasshoppers are being cultivated by commercial companies operating grasshopper farms and are being used as food and protein supplements.
As pests
Grasshoppers eat large quantities of foliage both as adults and during their development, and can be serious pests of arid land and prairies. Pasture, grain, forage, vegetable and other crops can be affected. Grasshoppers often bask in the sun, and thrive in warm sunny conditions, so drought stimulates an increase in grasshopper populations. A single season of drought is not normally sufficient to stimulate a major population increase, but several successive dry seasons can do so, especially if the intervening winters are mild so that large numbers of nymphs survive. Although sunny weather stimulates growth, there needs to be an adequate food supply for the increasing grasshopper population. This means that although precipitation is needed to stimulate plant growth, prolonged periods of cloudy weather will slow nymphal development.
Grasshoppers can best be prevented from becoming pests by manipulating their environment. Shade provided by trees will discourage them and they may be prevented from moving onto developing crops by removing coarse vegetation from fallow land and field margins and discouraging thick growth beside ditches and on roadside verges. With increasing numbers of grasshoppers, predator numbers may increase, but this seldom happens rapidly enough to have much effect on populations. Biological control is being investigated, and spores of the protozoan parasite Nosema locustae can be used mixed with bait to control grasshoppers, being more effective with immature insects. On a small scale, neem products can be effective as a feeding deterrent and as a disruptor of nymphal development. Insecticides can be used, but adult grasshoppers are difficult to kill, and as they move into fields from surrounding rank growth, crops may soon become reinfested.
Some grasshopper species, like the Chinese rice grasshopper, are a pest in rice paddies. Ploughing exposes the eggs on the surface of the field, to be destroyed by sunshine or eaten by natural enemies. Some eggs may be buried too deeply in the soil for hatching to take place.
Locust plagues can have devastating effects on human populations, causing famines and population upheavals. They are mentioned in both the Qur’an and the Bible and have also been held responsible for cholera epidemics, resulting from the corpses of locusts drowned in the Mediterranean Sea and decomposing on beaches. The FAO and other organisations monitor locust activity around the world. Timely application of pesticides can prevent nomadic bands of hoppers from forming before dense swarms of adults can build up. Besides conventional control using contact insecticides, biological pest control using the entomopathogenic fungus Metarhizium acridum, which specifically infects grasshoppers, has been used with some success.
Detection of explosives
In February 2020, researchers from Washington University in St. Louis announced they had engineered "cyborg grasshoppers" capable of accurately detecting explosives. In the project, funded by the US Office of Naval Research, researchers fitted grasshoppers with lightweight sensor backpacks that recorded and transmitted the electrical activity of their antennal lobes to a computer. According to the researchers, the grasshoppers were able to detect the location of the highest concentration of explosives. The researchers also tested the effect of combining sensorial information from several grasshoppers on detection accuracy. The neural activity from seven grasshoppers yielded an average detection accuracy rate of 80%, whereas a single grasshopper yielded a 60% rate.
In literature
The Egyptian word for locust or grasshopper was written snḥm in the consonantal hieroglyphic writing system. The pharaoh Ramesses II compared the armies of the Hittites to locusts: "They covered the mountains and valleys and were like locusts in their multitude."
One of Aesop's Fables, later retold by La Fontaine, is the tale of The Ant and the Grasshopper. The ant works hard all summer, while the grasshopper plays. In winter, the ant is ready but the grasshopper starves. Somerset Maugham's short story "The Ant and the Grasshopper" explores the fable's symbolism via complex framing. Other human weaknesses besides improvidence have become identified with the grasshopper's behaviour. So an unfaithful woman (hopping from man to man) is "a grasshopper" in "Poprygunya", an 1892 short story by Anton Chekhov, and in Jerry Paris's 1969 film The Grasshopper.
In mechanical engineering
The name "Grasshopper" was given to the Aeronca L-3 and Piper L-4 light aircraft, both used for reconnaissance and other support duties in World War II. The name is said to have originated when Major General Innis P. Swift saw a Piper making a rough landing and remarked that it looked like a grasshopper for its bouncing progress.
Grasshopper beam engines were beam engines pivoted at one end, the long horizontal arm resembling the hind leg of a grasshopper. The type was patented by William Freemantle in 1803.
| Biology and health sciences | Orthoptera | null |
10964778 | https://en.wikipedia.org/wiki/Cinder%20cone | Cinder cone | A cinder cone (or scoria cone) is a steep conical hill of loose pyroclastic fragments, such as volcanic clinkers, volcanic ash, or scoria that has been built around a volcanic vent. The pyroclastic fragments are formed by explosive eruptions or lava fountains from a single, typically cylindrical, vent. As the gas-charged lava is blown violently into the air, it breaks into small fragments that solidify and fall as either cinders, clinkers, or scoria around the vent to form a cone that often is symmetrical; with slopes between 30 and 40°; and a nearly circular ground plan. Most cinder cones have a bowl-shaped crater at the summit.
Mechanics of eruption
Cinder cones range in size from tens to hundreds of meters tall. They are composed of loose pyroclastic material (cinder or scoria), which distinguishes them from spatter cones, which are composed of agglomerated volcanic bombs.
The pyroclastic material making up a cinder cone is usually basaltic to andesitic in composition. It is often glassy and contains numerous gas bubbles "frozen" into place as magma exploded into the air and then cooled quickly. Lava fragments larger than 64 mm across, known as volcanic bombs, are also a common product of cinder cone eruptions.
The growth of a cinder cone may be divided into four stages. In the first stage, a low-rimmed scoria ring forms around the erupting event. During the second stage, the rim is built up and a talus slope begins to form outside the rim. The third stage is characterized by slumping and blasts that destroy the original rim, while the fourth stage is characterized by the buildup of talus beyond the zone where cinder falls to the surface (the ballistic zone).
During the waning stage of a cinder cone eruption, the magma has lost most of its gas content. This gas-depleted magma does not fountain but oozes quietly into the crater or beneath the base of the cone as lava. Lava rarely issues from the top (except as a fountain) because the loose, uncemented cinders are too weak to support the pressure exerted by molten rock as it rises toward the surface through the central vent. Because it contains so few gas bubbles, the molten lava is denser than the bubble-rich cinders. Thus, it often burrows out along the bottom of the cinder cone, lifting the less dense cinders like corks on water, and advances outward, creating a lava flow around the cone's base. When the eruption ends, a symmetrical cone of cinders sits at the center of a surrounding pad of lava. If the crater is fully breached, the remaining walls form an amphitheater or horseshoe shape around the vent.
Occurrence
Basaltic cinder cones are the most characteristic type of volcano associated with intraplate volcanism. They are particularly common in association with alkaline magmatism, in which the erupted lava is enriched in sodium and potassium oxides.
Cinder cones are also commonly found on the flanks of shield volcanoes, stratovolcanoes, and calderas. For example, geologists have identified nearly 100 cinder cones on the flanks of Mauna Kea, a shield volcano located on the island of Hawaii. Such cinder cones likely represent the final stages of activity of a mafic volcano. However, most volcanic cones formed in Hawaiian-type eruptions are spatter cones rather than cinder cones, due to the fluid nature of the lava.
The most famous cinder cone, Paricutin, grew out of a corn field in Mexico in 1943 from a new vent. Eruptions continued for nine years, built the cone to a height of , and produced lava flows that covered .
The Earth's most historically active cinder cone is Cerro Negro in Nicaragua. It is part of a group of four young cinder cones NW of Las Pilas volcano. Since its initial eruption in 1850, it has erupted more than 20 times, most recently in 1995 and 1999.
Satellite images suggest that cinder cones occur on other terrestrial bodies in the solar system. On Mars, they have been reported on the flanks of Pavonis Mons in Tharsis, in the region of Hydraotes Chaos on the bottom of the Coprates Chasma, or in the volcanic field Ulysses Colles. It is also suggested that domical structures in Marius Hills (on the Moon) might represent lunar cinder cones.
Effect of environmental conditions
The size and shape of cinder cones depend on environmental properties as different gravity and/or atmospheric pressure might change the dispersion of ejected scoria particles. For example, cinder cones on Mars seem to be more than two times wider than terrestrial analogues as lower atmospheric pressure and gravity enable wider dispersion of ejected particles over a larger area. Therefore, it seems that erupted amount of material is not sufficient on Mars for the flank slopes to attain the angle of repose and Martian cinder cones seem to be ruled mainly by ballistic distribution and not by material redistribution on flanks as typical on Earth.
Cinder cones often are highly symmetric, but strong prevailing winds at the time of eruption can cause a greater accumulation of cinder on the downwind side of the vent.
Monogenetic cones
Some cinder cones are monogenetic, forming from a single short eruptive episode that produces a very small volume of lava. The eruption typically last just weeks or months, but can occasionally last fifteen years or longer. Parícutin in Mexico, Diamond Head, Koko Head, Punchbowl Crater, Mt Le Brun from the Coalstoun Lakes volcanic field, and some cinder cones on Mauna Kea are monogenetic cinder cones. However, not all cinder cones are monogenetic, with some ancient cinder cones showing intervals of soil formation between flows that indicate that eruptions were separated by thousands to tens of thousands of years.
Monogenetic cones likely form when the rate of magma supply to a volcanic field is very low and the eruptions are spread out in space and time. This prevents any one eruption from establishing a system of "plumbing" that would provide an easy path to the surface for subsequent eruptions. Thus each eruption must find its independent path to the surface.
| Physical sciences | Volcanology | Earth science |
10970082 | https://en.wikipedia.org/wiki/Learning%20disability | Learning disability | Learning disability, learning disorder, or learning difficulty (British English) is a condition in the brain that causes difficulties comprehending or processing information and can be caused by several different factors. Given the "difficulty learning in a typical manner", this does not exclude the ability to learn in a different manner. Therefore, some people can be more accurately described as having a "learning difference", thus avoiding any misconception of being disabled with a possible lack of an ability to learn and possible negative stereotyping. In the United Kingdom, the term "learning disability" generally refers to an intellectual disability, while conditions such as dyslexia and dyspraxia are usually referred to as "learning difficulties".
While learning disability and learning disorder are often used interchangeably, they differ in many ways. Disorder refers to significant learning problems in an academic area. These problems, however, are not enough to warrant an official diagnosis. Learning disability, on the other hand, is an official clinical diagnosis, whereby the individual meets certain criteria, as determined by a professional (such as a psychologist, psychiatrist, speech-language pathologist, or paediatrician). The difference is in the degree, frequency, and intensity of reported symptoms and problems, and thus the two should not be confused. When the term "learning disorder" is used, it describes a group of disorders characterized by inadequate development of specific academic, language, and speech skills. Types of learning disorders include reading (dyslexia), arithmetic (dyscalculia) and writing (dysgraphia).
The unknown factor is the disorder that affects the brain's ability to receive and process information. This disorder can make it problematic for a person to learn as quickly or in the same way as someone who is not affected by a learning disability. People with a learning disability have trouble performing specific types of skills or completing tasks if left to figure things out by themselves or if taught in conventional ways.
Individuals with learning disabilities can face unique challenges that are often pervasive throughout the lifespan. Depending on the type and severity of the disability, interventions, and current technologies may be used to help the individual learn strategies that will foster future success. Some interventions can be quite simple, while others are intricate and complex. Current technologies may require student training to be effective classroom supports. Teachers, parents, and schools can create plans together that tailor intervention and accommodations to aid the individuals in successfully becoming independent learners. A multi-disciplinary team frequently helps to design the intervention and to coordinate the execution of the intervention with teachers and parents. This team frequently includes school psychologists, special educators, speech therapists (pathologists), occupational therapists, psychologists, ESL teachers, literacy coaches, and/or reading specialists.
Definition
In the United States, a committee of representatives of organizations committed to the education and welfare of individuals with learning disabilities is known as the National Joint Committee on Learning Disabilities (NJCLD). The NJCLD used the term 'learning disability' to indicate a discrepancy between a child's apparent capacity to learn and their level of achievement. Several difficulties existed, however, with the NJCLD standard of defining learning disability. One such difficulty was its belief of central nervous system dysfunction as a basis of understanding and diagnosing learning disability. This conflicted with the fact that many individuals who experienced central nervous system dysfunction, such as those with cerebral palsy, did not experience disabilities in learning. On the other hand, those individuals who experienced multiple handicapping conditions along with learning disability frequently received inappropriate assessment, planning, and instruction. The NJCLD notes that it is possible for learning disability to occur simultaneously with other handicapping conditions, however, the two should not be directly linked together or confused.
In the 1980s, NJCLD, therefore, defined the term learning disability as:
The 2002 LD Roundtable produced the following definition:The issue of defining learning disabilities has generated significant and ongoing controversy. The term "learning disability" does not exist in DSM-IV, but it has been added to the DSM-5. The DSM-5 does not limit learning disorders to a particular diagnosis such as reading, mathematics, or written expression. Instead, it is a single diagnosis criterion describing drawbacks in general academic skills and includes detailed specifiers for the areas of reading, mathematics, and written expression.
United States and Canada
In the United States and Canada, the terms learning disability and learning disorder (LD) refer to a group of disorders that affect a broad range of academic and functional skills including the ability to speak, listen, read, write, spell, reason, organize information, and do math. People with learning disabilities generally have average or higher intelligence.
Legislation in the United States
The Section 504 of the Rehabilitation Act 1973, effective May 1977, guarantees certain rights to people with disabilities, especially in the cases of education and work, such being in schools, colleges and university settings.
The Individuals with Disabilities Education Act, formerly known as the Education for All Handicapped Children Act, is a United States federal law that governs how states and public agencies provide early intervention, special education and related services to children with disabilities. It addresses the educational needs of children with disabilities from birth to the age of 21. Considered as a civil rights law, states are not required to participate.
Policymaking in Canada
In Canada, the first association in support of children with learning disabilities was founded in 1962 by a group of concerned parents. Originally called the Association for Children with Learning Disabilities, the Learning Disabilities Association of Canada – LDAC was created to provide awareness and services for individuals with learning disabilities, their families, at work, and the community. Since education is largely the responsibility of each province and territory in Canada, provinces and territories have jurisdiction over the education of individuals with learning disabilities, which allows the development of policies and support programs that reflect the unique multicultural, linguistic, and socioeconomic conditions of its area.
United Kingdom
In the UK, terms such as specific learning difficulty (SpLD), developmental dyslexia, developmental coordination disorder and dyscalculia are used to cover the range of learning difficulties referred to in the United States as "learning disabilities". In the UK, the term "learning disability" refers to a range of developmental disabilities or conditions that are almost invariably associated with more severe generalized cognitive impairment. The Lancet defines 'learning disability' as a "significant general impairment in intellectual functioning acquired during childhood", and states that roughly one in 50 British adults have one.
Japan
In Japan, acknowledgement and support for students with learning disabilities has been a fairly recent development, and has improved drastically since the start of the 21st century. The first definition for learning disability was coined in 1999, and in 2001, the Enrichment Project for the Support System for Students with Learning Disabilities was established. Since then, there have been significant efforts to screen children for learning disabilities, provide follow-up support, and provide networking between schools and specialists.
Effects
The effects of having a learning disability or learning difference are not limited to educational outcomes: individuals with learning disabilities may experience social problems as well. Neuropsychological differences can affect the accurate perception of social cues with peers. Researchers argue persons with learning disabilities not only experience negative effects as a result of their learning distinctions, but also as a result of carrying a stigmatizing label. It has generally been difficult to determine the efficacy of special education services because of data and methodological limitations. Emerging research suggests adolescents with learning disabilities experience poorer academic outcomes even compared to peers who began high school with similar levels of achievement and comparable behaviors. It seems their poorer outcomes may be at least partially due to the lower expectations of their teachers; national data show teachers hold expectations for students labeled with learning disabilities that are inconsistent with their academic potential (as evidenced by test scores and learning behaviors). It has been said that there is a strong connection between children with a learning disability and their educational performance.
Many studies have been done to assess the correlation between learning disability and self-esteem. These studies have shown that an individual's self-esteem is indeed affected by their own awareness of their learning disability. Students with a positive perception of their academic abilities generally tend to have higher self-esteem than those who do not, regardless of their actual academic achievement. However, studies have also shown that several other factors can influence self-esteem. Skills in non-academic areas, such as athletics and arts, improve self-esteem. Also, a positive perception of one's physical appearance has also been shown to have positive effects of self-esteem. Another important finding is that students with learning disabilities are able to distinguish between academic skill and intellectual capacity. This demonstrates that students who acknowledge their academic limitations but are also aware of their potential to succeed in other intellectual tasks see themselves as intellectually competent individuals, which increases their self-esteem.
Research involving individuals with learning disabilities who exhibit challenging behaviors who are subsequently treated with antipsychotic medications provides little evidence that any benefits outweigh the risk.
Causes
The causes for learning disabilities are not well understood, and sometimes there is no apparent cause for a learning disability. However, some causes of neurological impairments include:
Heredity and genetics: Learning disabilities are often linked through genetics and run in the family. Children who have learning disabilities often have parents who have the same struggles. Children of parents who had less than 12 years of school are more likely to have a reading disability. Some children have spontaneous mutations (i.e. not present in either parent) which can cause developmental disorders including learning disabilities. One study estimated that about one in 300 children had such spontaneous mutations, for example a fault in the CDK13 gene which is associated with learning and communication difficulties in the children affected.
Problems during pregnancy and birth: A learning disability can result from anomalies in the developing brain, illness or injury. Risk factors are foetal exposure to alcohol or drugs and low birth weight (3 pounds or less). These children are more likely to develop a disability in math or reading. Children who are born prematurely, late, have a longer labor than usual, or have trouble receiving oxygen are more likely to develop a learning disability.
Accidents after birth: Learning disabilities can also be caused by head injuries, malnutrition, or by toxic exposure (such as heavy metals or pesticides).
Diagnosis
IQ-achievement discrepancy
Learning disabilities can be identified by psychiatrists, speech language pathologists, school psychologists, clinical psychologists, counseling psychologists, neuropsychologists, speech language pathologists, and other learning disability specialists through a combination of intelligence testing, academic achievement testing, classroom performance, and social interaction and aptitude. Other areas of assessment may include perception, cognition, memory, attention, and language abilities. The resulting information is used to determine whether a child's academic performance is commensurate with their cognitive ability. If a child's cognitive ability is much higher than their academic performance, the student is often diagnosed with a learning disability. The DSM-IV and many school systems and government programs diagnose learning disabilities in this way (DSM-IV uses the term "disorder" rather than "disability").
Although the discrepancy model has dominated the school system for many years, there has been substantial criticism of this approach among researchers. Recent research has provided little evidence that a discrepancy between formally measured IQ and achievement is a clear indicator of LD. Furthermore, diagnosing on the basis of a discrepancy does not predict the effectiveness of treatment. Low academic achievers who do not have a discrepancy with IQ (i.e. their IQ scores are also low) appear to benefit from treatment just as much as low academic achievers who do have a discrepancy with IQ (i.e. their IQ scores are higher than their academic performance would suggest).
Since 1998, there have been attempts to create a reference index more useful than IQ to generate predicted scores on achievement tests. For example, for a student whose vocabulary and general knowledge scores matches their reading comprehension score a teacher could assume that reading comprehension can be supported through work in vocabulary and general knowledge. If the reading comprehension score is lower in the appropriate statistical sense it would be necessary to first rule out things like vision problems.
Response to intervention
Much current research has focused on a treatment-oriented diagnostic process known as response to intervention (RTI). Researcher recommendations for implementing such a model include early screening for all students, placing those students who are having difficulty into research-based early intervention programs, rather than waiting until they meet diagnostic criteria. Their performance can be closely monitored to determine whether increasingly intense intervention results in adequate progress. Those who respond will not require further intervention. Those who do not respond adequately to regular classroom instruction (often called "Tier 1 instruction") and a more intensive intervention (often called "Tier 2" intervention) are considered "non-responders." These students can then be referred for further assistance through special education, in which case they are often identified with a learning disability. Some models of RTI include a third tier of intervention before a child is identified as having a learning disability.
A primary benefit of such a model is that it would not be necessary to wait for a child to be sufficiently far behind to qualify for assistance. This may enable more children to receive assistance before experiencing significant failure, which may, in turn, result in fewer children who need intensive and expensive special education services. In the United States, the 2004 reauthorization of the Individuals with Disabilities Education Act permitted states and school districts to use RTI as a method of identifying students with learning disabilities. RTI is now the primary means of identification of learning disabilities in Florida.
The process does not take into account children's individual neuropsychological factors such as phonological awareness and memory, that can inform design instruction. By not taking into account specific cognitive processes, RTI fails to inform educators about a students' relative strengths and weaknesses Second, RTI by design takes considerably longer than established techniques, often many months to find an appropriate tier of intervention. Third, it requires a strong intervention program before students can be identified with a learning disability. Lastly, RTI is considered a regular education initiative and consists of members of general education teachers, in conjunction with other qualified professionals. Occupational therapists in particular can support students in the educational setting by helping children in academic and non-academic areas of school including the classroom, recess and meal time. They can provide strategies, therapeutic interventions, suggestions for adaptive equipment, and environmental modifications. Occupational therapists can work closely with the child's teacher and parents to facilitate educational goals specific to each child under an RTI and/or IEP.
Latino English language learners
Demographers in the United States report that there has been a significant increase in immigrant children in the United States over the past two decades. This information is vital because it has been and will continue to affect both students and how educators approach teaching methods. Various teaching strategies are more successful for students that are linguistic or culturally diverse versus traditional methods of teaching used for students whose first language is English. It is then also true that the proper way to diagnose a learning disability in English language learners (ELL) differs. In the United States, there has been a growing need to develop the knowledge and skills necessary to provide effective school psychological services, specifically for those professionals who work with immigrant populations.
Currently, there are no standardized guidelines for the process of diagnosing ELL with specific learning disabilities (SLD). This is a problem since many students will fall through the cracks as educators are unable to clearly assess if a student's delay is due to a language barrier or true learning disability. Without a clear diagnosis, many students will suffer because they will not be provided with the tools they need to succeed in the public education school system. For example, in many occasions teachers have suggested retention or have taken no action at all when they lack experience working with English language learners. Students were commonly pushed toward testing, based on an assumption that their poor academic performance or behavioral difficulties indicated a need for special education. Linguistically responsive psychologist understand that second language acquisition is a process and they understand how to support ELLs' growth in language and academically. When ELLs are referred for a psychoeducational assessment, it is difficult to isolate and disentangle what are the effects of the language acquisition process, from poor quality educational services, from what may be academic difficulties that result from processing disorders, attention problems, and learning disabilities. Additionally not having trained staff and faculty becomes more of an issue when staff is unaware of numerous types of psychological factors that immigrant children in the U.S. could be potentially dealing with. These factors that include acculturation, fear and/or worry of deportation, separation from social supports such as parents, language barriers, disruptions in learning experiences, stigmatization, economic challenge, and risk factors associated with poverty. In the United States, there are no set policies mandating that all districts employ bilingual school psychologist, nor are schools equipped with specific tools and resources to assist immigrant children and families. Many school districts do not have the proper personnel that is able to communicate with this population.
Spanish-speaking ELL
A well trained bilingual school psychologist will be able to administer and interpret assessments using psychological testing tools. Also, an emphasis is placed on informal assessment measures such as language samples, observations, interviews, and rating scales as well as curriculum-based measurement to complement information gathered from formal assessments. A compilation of these tests is used to assess whether an ELL student has a learning disability or merely is academically delayed because of language barriers or environmental factors. Many schools do not have a school psychologist with the proper training nor access to appropriate tools. Also, many school districts frown upon taking the appropriate steps to diagnosing ELL students.
Assessment
Many normed assessments can be used in evaluating skills in the primary academic domains: reading, including word recognition, fluency, and comprehension; mathematics, including computation and problem solving; and written expression, including handwriting, spelling and composition.
The most commonly used comprehensive achievement tests include the Woodcock-Johnson IV (WJ IV), Wechsler Individual Achievement Test II (WIAT II), the Wide Range Achievement Test III (WRAT III), and the Stanford Achievement Test–10th edition. These tests include measures of many academic domains that are reliable in identifying areas of difficulty.
In the reading domain, there are also specialized tests that can be used to obtain details about specific reading deficits. Assessments that measure multiple domains of reading include Gray's Diagnostic Reading Tests–2nd edition (GDRT II) and the Stanford Diagnostic Reading Assessment. Assessments that measure reading subskills include the Gray Oral Reading Test IV – Fourth Edition (GORT IV), Gray Silent Reading Test, Comprehensive Test of Phonological Processing (CTOPP), Tests of Oral Reading and Comprehension Skills (TORCS), Test of Reading Comprehension 3 (TORC-3), Test of Word Reading Efficiency (TOWRE), and the Test of Reading Fluency. A more comprehensive list of reading assessments may be obtained from the Southwest Educational Development Laboratory.
The purpose of assessment is to determine what is needed for intervention, which also requires consideration of contextual variables and whether there are comorbid disorders that must also be identified and treated, such as behavioral issues or language delays. These contextual variables are often assessed using parent and teacher questionnaire forms that rate the students' behaviors and compares them to standardized norms.
However, caution should be made when suspecting the person with a learning disability may also have dementia, especially as people with Down's syndrome may have the neuroanatomical profile but not the associated clinical signs and symptoms. Examination can be carried out of executive functioning as well as social and cognitive abilities but may need adaptation of standardized tests to take account of special needs.
Types
Learning disabilities can be categorized by either the type of information processing affected by the disability or by the specific difficulties caused by a processing deficit.
By stage of information processing
Learning disabilities fall into broad categories based on the four stages of information processing used in learning: input, integration, storage, and output. Many learning disabilities are a compilation of a few types of abnormalities occurring at the same time, as well as with social difficulties and emotional or behavioral disorders.
Input: This is the information perceived through the senses, such as visual and auditory perception. Difficulties with visual perception can cause problems with recognizing the shape, position, or size of items seen. There can be problems with sequencing, which can relate to deficits with processing time intervals or temporal perception. Difficulties with auditory perception can make it difficult to screen out competing sounds in order to focus on one of them, such as the sound of the teacher's voice in a classroom setting. Some children appear to be unable to process tactile input. For example, they may seem insensitive to pain or dislike being touched.
Integration: This is the stage during which perceived input is interpreted, categorized, placed in a sequence, or related to previous learning. Students with problems in these areas may be unable to tell a story in the correct sequence, unable to memorize sequences of information such as the days of the week, able to understand a new concept but be unable to generalize it to other areas of learning, or able to learn facts but be unable to put the facts together to see the "big picture." A poor vocabulary may contribute to problems with comprehension.
Storage: Problems with memory can occur with short-term or working memory, or with long-term memory. Most memory difficulties occur with one's short-term memory, which can make it difficult to learn new material without more repetitions than usual. Difficulties with visual memory can impede learning to spell.
Output: Information comes out of the brain either through words, that is, language output, or through muscle activity, such as gesturing, writing or drawing. Difficulties with language output can create problems with spoken language. Such difficulties include answering a question on demand, in which one must retrieve information from storage, organize our thoughts, and put the thoughts into words before we speak. It can also cause trouble with written language for the same reasons. Difficulties with motor abilities can cause problems with gross and fine motor skills. People with gross motor difficulties may be clumsy, that is, they may be prone to stumbling, falling, or bumping into things. They may also have trouble running, climbing, or learning to ride a bicycle. People with fine motor difficulties may have trouble with handwriting, buttoning shirts, or tying shoelaces.
By function impaired
Deficits in any area of information processing can manifest in a variety of specific learning disabilities (SLD). It is possible for an individual to have more than one of these difficulties. This is referred to as comorbidity or co-occurrence of learning disabilities. In the UK, the term dual diagnosis is often used to refer to co-occurrence of learning difficulties.
Reading disorder (ICD-10 and DSM-IV codes: F81.0/315.00)
Reading disorder is the most common learning disability. Of all students with specific learning disabilities, 70–80% have deficits in reading. The term "Developmental Dyslexia" is often used as a synonym for reading disability; however, many researchers assert that there are different types of reading disabilities, of which dyslexia is one. A reading disability can affect any part of the reading process, including difficulty with accurate or fluent word recognition, or both, word decoding, reading rate, prosody (oral reading with expression), and reading comprehension. Before the term "dyslexia" came to prominence, this learning disability used to be known as "word blindness."
Common indicators of reading disability include difficulty with phonemic awareness—the ability to break up words into their component sounds, and difficulty with matching letter combinations to specific sounds (sound-symbol correspondence).
Disorder of written expression (ICD-10 and DSM-IV-TR codes 315.2)
The DSM-IV-TR criteria for a disorder of written expression is writing skills (as measured by a standardized test or functional assessment) that fall substantially below those expected based on the individual's chronological age, measured intelligence, and age-appropriate education, (Criterion A). This difficulty must also cause significant impairment to academic achievement and tasks that require composition of written text (Criterion B), and if a sensory deficit is present, the difficulties with writing skills must exceed those typically associated with the sensory deficit, (Criterion C).
Individuals with a diagnosis of a disorder of written expression typically have a combination of difficulties in their abilities with written expression as evidenced by grammatical and punctuation errors within sentences, poor paragraph organization, multiple spelling errors, and excessively poor penmanship. A disorder in spelling or handwriting without other difficulties of written expression do not generally qualify for this diagnosis. If poor handwriting is due to an impairment in the individuals' motor coordination, a diagnosis of developmental coordination disorder should be considered.
By a number of organizations, the term "dysgraphia" has been used as an overarching term for all disorders of written expression.
Math disability (ICD-10 and DSM-IV codes F81.2-3/315.1)
Sometimes called dyscalculia, a math disability involves difficulties such as learning math concepts (such as quantity, place value, and time), difficulty memorizing math facts, difficulty organizing numbers, and understanding how problems are organized on the page. Dyscalculics are often referred to as having poor "number sense".
Non ICD-10/DSM
Nonverbal learning disability: Nonverbal learning disabilities often manifest in motor clumsiness, poor visual-spatial skills, problematic social relationships, difficulty with mathematics, and poor organizational skills. These individuals often have specific strengths in the verbal domains, including early speech, large vocabulary, early reading and spelling skills, excellent rote memory and auditory retention, and eloquent self-expression.
Disorders of speaking and listening: Difficulties that often co-occur with learning disabilities include difficulty with memory, social skills and executive functions (such as organizational skills and time management).
Management
Interventions include:
Mastery model:
Learners work at their own level of mastery.
Practice
Gain fundamental skills before moving onto the next level
Note: this approach is most likely to be used with adult learners or outside the mainstream school system.
Direct instruction:
Emphasizes carefully planned lessons for small learning increments
Scripted lesson plans
Rapid-paced interaction between teacher and students
Correcting mistakes immediately
Achievement-based grouping
Frequent progress assessments
Classroom adjustments:
Special seating assignments
Alternative or modified assignments
Modified testing procedures
Quiet environment
Special equipment:
Word processors with spell checkers and dictionaries
Text-to-speech and speech-to-text programs
Talking calculators
Books on tape
Computer-based activities
Classroom assistants:
Note-takers
Readers
Proofreaders
Scribes
Special education:
Prescribed hours in a resource room
Placement in a resource room
Enrollment in a special school or a separate classroom in a regular school for learning disabled students
Individual education plan (IEP)
Educational therapy
Sternberg has argued that early remediation can greatly reduce the number of children meeting diagnostic criteria for learning disabilities. He has also suggested that the focus on learning disabilities and the provision of accommodations in school fails to acknowledge that people have a range of strengths and weaknesses, and places undue emphasis on academic success by insisting that people should receive additional support in this arena but not in music or sports. Other research has pinpointed the use of resource rooms as an important—yet often politicized component of educating students with learning disabilities.
Helping individuals with learning disabilities
Many individuals with learning disabilities may not openly disclose their condition. Some experts say that an instructor directly asking or assuming potential disabilities could cause potential harm to an individual's self esteem. In addition, if information about certain disabilities were made aware, it may be beneficial to be mindful about one's approach regarding the disability and avoid vocabulary that may insinuate that the learning disability is an obstacle or shortcoming as this may potentially be harmful to an individual's mental health and self esteem. Research suggests that accumulating positive experiences such as success in interpersonal relationships, achievements, and overcoming stress leads to the formation of self-esteem leading to the acceptance of one's disability and a better life outcome. This suggests that working with the disability may result in more positive outcomes rather than attempting to fix it. As an instructor or tutor, it may be helpful to consider asking the needs of individuals with disabilities as they know their disability the best. Some question to consider:
What part of the assignment do you want to focus on?
Where in our space would you most prefer to work?
What tools or technologies do you tend to use most frequently when you write?
Are you comfortable reading your paper out loud or would you prefer if I read it?
How do you learn best (i.e. Do you learn best by doing, seeing, or hearing)?
Society and culture (United States)
School laws
Schools in the United States have a legal obligation to new arrivals to the country, including undocumented students. The landmark Supreme Court ruling Plyler v. Doe (1982) grants all children, no matter their legal status, the right to a free education. Additionally, specifically in regards to ELLs, the supreme court ruling Lau v. Nichols (1974) stated that equal treatment in school did not mean equal educational opportunity. This ruling is also supported by English language development services provided in schools, but these rulings do not require the individuals that teach and provide services to have any specific training nor is licensing different from a typical teacher or services provider.
Issues Regarding Standardized Testing
Problems still exist regarding the fairness of standardized testing. Providing testing accommodations to students with learning disabilities has become increasingly common. One of such issues that introduce iniquity to those with disabilities is the handwriting bias. The handwriting bias involves the tendency of raters to identify more personally with authors of handwritten essays compared to word-processed essays resulting in awarding a higher rating to the handwritten essays despite both essays being identical in terms of content. Several studies have analyzed the differences in standardized scores of handwriting and word-processed (typed) essays between students with and without disability. Results suggest handwritten essays of students with and without disabilities consistently received higher scores compared to word processed versions.
Critique of the medical model
Learning disability theory is founded in the medical model of disability, in that disability is perceived as an individual deficit that is biological in origin. Researchers working within a social model of disability assert that there are social or structural causes of disability or the assignation of the label of disability, and even that disability is entirely socially constructed. Since the turn of the 19th century, education in the United States has been geared toward producing citizens who can effectively contribute to a capitalistic society, with a cultural premium on efficiency and science. More agrarian cultures, for example, do not even use learning ability as a measure of adult adequacy, whereas the diagnosis of learning disabilities is prevalent in Western capitalistic societies because of the high value placed on speed, literacy, and numeracy in both the labor force and school system.
Culture
There are three patterns that are well known in regards to mainstream students and minority labels in the United States:
"A higher percentage of minority children than of white children are assigned to special education";
"within special education, white children are assigned to less restrictive programs than are their minority counterparts";
"the data — driven by inconsistent methods of diagnosis, treatment, and funding — make the overall system difficult to describe or change".
In the present day, it has been reported that white districts have more children from minority backgrounds enrolled in special education than they do majority students. "It was also suggested that districts with a higher percentage of minority faculty had fewer minority students placed in special education suggesting that 'minority students are treated differently in predominantly white districts than in predominantly minority districts'".
Educators have only recently started to look into the effects of culture on learning disabilities. If a teacher ignores a student's culturally diverse background, the student will suffer in the class. "The cultural repertoires of students from cultural learning disorder backgrounds have an impact on their learning, school progress, and behavior in the classroom". These students may then act out and not excel in the classroom and will, therefore, be misdiagnosed: "Overall, the data indicates that there is a persistent concern regarding the misdiagnosis and inappropriate placement of students from diverse backgrounds in special education classes since the 1975".
Social roots of learning disabilities in the U.S.
Learning disabilities have a disproportionate identification of racial and ethnic minorities and students who have low socioeconomic status (SES). While some attribute the disproportionate identification of racial/ethnic minorities to racist practices or cultural misunderstanding, others have argued that racial/ethnic minorities are overidentified because of their lower status. Similarities were noted between the behaviors of "brain-injured" and lower class students as early as the 1960s. The distinction between race/ethnicity and SES is important to the extent that these considerations contribute to the provision of services to children in need. While many studies have considered only one characteristic of the student at a time, or used district- or school-level data to examine this issue, more recent studies have used large national student-level datasets and sophisticated methodology to find that the disproportionate identification of African American students with learning disabilities can be attributed to their average lower SES, while the disproportionate identification of Latino youth seems to be attributable to difficulties in distinguishing between linguistic proficiency and learning ability. Although the contributing factors are complicated and interrelated, it is possible to discern which factors really drive disproportionate identification by considering a multitude of student characteristics simultaneously. For instance, if high SES minorities have rates of identification that are similar to the rates among high SES Whites, and low SES minorities have rates of identification that are similar to the rates among low SES Whites, we can know that the seemingly higher rates of identification among minorities result from their greater likelihood to have low SES. Summarily, because the risk of identification for White students who have low SES is similar to that of Black students who have low SES, future research and policy reform should focus on identifying the shared qualities or experiences of low SES youth that lead to their disproportionate identification, rather than focusing exclusively on racial/ethnic minorities. It remains to be determined why lower SES youth are at higher risk of incidence, or possibly just of identification, with learning disabilities.
Learning disabilities in adulthood
A common misconception about those with learning disabilities is that they outgrow it as they enter adulthood. This is often not the case and most adults with learning disabilities still require resources and care to help manage their disability. One resource available is the Adult Basic Education (ABE) programs, at the state level. ABE programs are allotted certain amounts of funds per state in order to provide resources for adults with learning disabilities. This includes resources to help them learn basic life skills in order to provide for themselves. ABE programs also provide help for adults who lack a high school diploma or an equivalent. These programs teach skills to help adults get into the workforce or into a further level of education. There is a certain pathway that these adults and instructors should follow in order to ensure these adults have the abilities needed to succeed in life. Some ABE programs offer GED preparation programs to support adults through the process to get a GED. It is important to note that ABE programs do not always have the expected outcome on things like employment. Participants in ABE programs are given tools to help them succeed and get a job but, employment is dependent on more than just a guarantee of a job post-ABE. Employment varies based on the level of growth a participant experiences in an ABE program, the personality and behavior of the participant, and the job market they are entering into following completion of an ABE program.
Another program to assist adults with disabilities are federal programs called "home and community based services" (HCBS). Medicaid funds these programs for many people through a fee waiver system, however, there are still lots of people on a stand-by list. These programs are primarily used for adults with Autism Spectrum Disorders. HCBS programs offer service more dedicated to caring for the adult, not so much providing resources for them to transition into the workforce. Some services provided are: therapy, social skills training, support groups, and counseling.
Contrast with other conditions
People with an IQ lower than 70 are usually characterized as having an intellectual disability and are not included under most definitions of learning disabilities because their difficulty in learning are considered to be related directly to their overall low intelligence.
Attention-deficit hyperactivity disorder (ADHD) is often studied in connection with learning disabilities, but it is not actually included in the standard definitions of learning disabilities. Individuals with ADHD may struggle with learning, but can often learn adequately once successfully treated for the ADHD. A person can have ADHD but not learning disabilities or have learning disabilities without having ADHD. The conditions can co-occur.
People diagnosed with ADHD sometimes have impaired learning. Some of the struggles people with ADHD have might include lack of motivation, high levels of anxiety, and the inability to process information. There are studies that suggest people with ADHD generally have a positive attitude toward academics and, with medication and developed study skills, can perform just as well as individuals without learning disabilities. Also, using alternate sources of gathering information, such as websites, study groups and learning centers, can help a person with ADHD be academically successful.
Before the discovery of ADHD, it was technically included in the definition of LDs since it has a very pronounced effect on the "executive functions" required for learning. Thus historically, ADHD was not clearly distinguished from other disabilities related to learning. Therefore, when a person presents with difficulties in learning, ADHD should be considered as well. Scientific research continues to explore the traits, struggles, effective learning styles and comorbid LDs of those with ADHD.
Learning disabilities affect the writing process
The ability to express one's thoughts and opinions in an organized fashion and in written form is an essential life skill that individuals have been taught and practiced repetitively since youth. The writing process includes, but is not limited to: understanding the genre, style, reading, critical thinking, writing and proofreading. In the case of individuals possessing a learning disability, deficits may be present that could impair the individuals' ability to carry out these necessary steps and express their thoughts in an organized manner. Reading is a crucial step to quality writing, oftentimes, it is practiced from a young age. Reading increases the attention span, allows exposure to a variety of genres and writing styles, and allows for the accumulation of a wide range of vocabulary.
Studies suggest that students with learning disabilities typically have difficulty with word recognition, the process of connecting the text to its meaning. This makes the reading process slow and cognitively laborious, which can be a very frustrating experience, causing students with learning disabilities to spend less time reading compared to their classmates. This in turn can negatively affect vocabulary acquisition and comprehension development of the individual.
In the context of standardized test taking, studies show that the strongest predictor of the level of performance during standardized essay writing was vocabulary complexity, specifically, the number of words with more than two syllables. Studies have suggested that individuals with ADHD tend to use simple structures and vocabulary. This puts many students with learning disabilities at a disadvantage since their knowledge of complex vocabulary usually does not compare to their peers. Based on such patterns, early interventions such as reading and writing curriculums from a young age could provide opportunities for vocabulary acquisition and development.
In addition, some students with learning disabilities tend to have difficulty separating the different stages of writing and devote little time to the planning stage. Oftentimes, they attempt to simultaneously reflect on their spelling while putting ideas together causing them to overload their attention system and make a number of spelling mistakes.
All together, the tendency of students with learning disabilities to dedicate little time to the planning and revision process compared to their peers often results in a lower level of coherence and quality of their written composition and a lower quality rating in the case for standardized tests. There is a lack of research in this area due to the complex relationship between the brain and one's ability to articulate ideas in writing. More research should be conducted in order to assess these factors and test the effectiveness of various intervention techniques.
| Biology and health sciences | Disabilities | Health |
341640 | https://en.wikipedia.org/wiki/Coriander | Coriander | Coriander (; Coriandrum sativum), also known as cilantro (), is an annual herb in the family Apiaceae.
Most people perceive coriander to have a tart, slightly citrus taste. Due to variations in the gene OR6A2, some people perceive it to have a soap-like taste, or even a pungent or rotten taste.
It is native to the Mediterranean Basin. All parts of the plant are edible, but the fresh leaves and the dried seeds are the parts most traditionally used in cooking.
Description
It is a soft plant growing to tall. The leaves are variable in shape, broadly lobed at the base of the plant, and slender and feathery higher on the flowering stems.
The flowers are borne in small umbels, white or very pale pink, asymmetrical, with the petals pointing away from the centre of the umbel longer () than those pointing toward it (only long). The fruit is a globular, dry schizocarp in diameter. The pollen size is approximately .
Taste and smell
The essential oil from coriander leaves and seeds contains mixed polyphenols and terpenes, including linalool as the major constituent accounting for the aroma and flavour of coriander.
Different people may perceive the taste of coriander leaves differently. Those who enjoy it say it has a refreshing, lemony or lime-like flavour, while those who dislike it have a strong aversion to its pungent taste and smell, characterizing it as soapy or rotten. Studies also show variations in preference among different ethnic groups: 21% of East Asians, 17% of Caucasians, and 14% of people of African descent expressed a dislike for coriander, but among the groups where coriander is popular in their cuisine, only 7% of South Asians, 4% of Hispanics, and 3% of Middle Eastern subjects expressed a dislike.
About 80% of identical twins shared the same preference for the herb, but fraternal twins agreed only about half the time, strongly suggesting a genetic component to the preference. In a genetic survey of nearly 30,000 people, two genetic variants linked to the perception of coriander have been found, the most common of which is a gene involved in sensing smells. The gene OR6A2 lies within a cluster of olfactory-receptor genes, and encodes a receptor that is highly sensitive to aldehyde chemicals. Flavour chemists have found that the coriander aroma is created by a half-dozen substances, most of which are aldehydes. Those who dislike the taste are sensitive to the offending unsaturated aldehydes and, at the same time, may be unable to detect the aromatic chemicals that others find pleasant. Association between its taste and several other genes, including a bitter-taste receptor, have also been found.
Similar plants
Eryngium foetidum, also a member of the Apiaceae, has a similar but more intense taste. Known as culantro and ngò gai, it is found in Mexico, the Caribbean, Central and South America, and South East Asia cuisine.
Persicaria odorata is commonly called Vietnamese coriander, or rau răm. The leaves have a similar odour and flavour to coriander. It is a member of the Polygonaceae, or buckwheat family.
Papaloquelite is one common name for Porophyllum ruderale subsp. macrocephalum, a member of the Asteraceae, the sunflower family. This species is found growing wild from Texas to Argentina.
Etymology
First attested in English during the late 14th century, the word "coriander" derives from the Old French , which comes from Latin , in turn from Ancient Greek (or ), possibly derived from or related to (a bed bug), and was given on account of its fetid, bug-like smell.
The earliest attested form of the word is the Mycenaean Greek (variants: , , ) written in Linear B syllabic script (reconstructed as , similar to the name of Minos' daughter Ariadne) which later evolved to koriannon or koriandron, and (German).
is the Spanish word for coriander, also deriving from coriandrum. It is the common term in US English for coriander leaves due to their extensive use in Mexican cuisine, but the seeds are referred to as coriander in American English.
Origin
Coriander grows wild over a wide area of Western Asia and Southern Europe, making it difficult to define where the plant is native and where it was only recently established. Recent works suggest that wild coriander in Israel and Portugal might be an ancestor of cultivated coriander. They have low germination rates and a small vegetative appearance. Israeli coriander has an extremely hard fruit coat.
In Israel, fifteen desiccated mericarps were found in the Pre-Pottery Neolithic B level (six to eight thousand years ago) of the Nahal Hemar Cave, and eleven from ~8,000–7,500 years ago in Pre-Pottery Neolithic C in Atlit-Yam. If these finds do belong to these archaeological layers, they are the oldest find of coriander in the world.
About of coriander mericarps were recovered from the tomb of Tutankhamen. As coriander does not grow wild in Egypt, this could be proof that coriander was cultivated by the ancient Egyptians. The Ebers Papyrus, an Egyptian text dated around 1550 BCE, mentioned uses of coriander.
Coriander may have been cultivated in Greece since at least the second millennium BCE. One of the Linear B tablets recovered from Pylos refers to the species as being cultivated for the manufacture of perfumes. It was used in two forms: as a spice for its seeds and as an herb for the flavour of its leaves. This appears to be confirmed by archaeological evidence: the large quantities of coriander retrieved from an Early Bronze Age layer at Sitagroi in Macedonia could point to cultivation of the herb at that time.
Allergies
Some people are allergic to coriander leaves or seeds, having symptoms similar to those of other food allergies. A cross-sectional study of 589 cases where food allergies to spices were suspected found 32% of pin-prick tests in children and 23% in adults were positive for coriander and other members of the family Apiaceae, including caraway, fennel, and celery. The allergic symptoms may be minor or life-threatening.
Uses
Nutrition
Raw coriander leaves are 92% water, 4% carbohydrates, 2% protein, and less than 1% fat. The nutritional profile of coriander seeds is different from that of fresh stems or leaves. In a reference amount, leaves are particularly rich in vitamin A, vitamin C, and vitamin K, with moderate content of dietary minerals. Although seeds generally have lower vitamin content, they do provide significant amounts of dietary fiber, calcium, selenium, iron, magnesium, and manganese.
Culinary
All parts of the plant are edible. Fresh leaves and dried seeds are the most commonly used in cooking. Coriander roots are an important element of Thai cooking. Coriander is used in cuisines throughout the world.
Leaves
The leaves are variously referred to as coriander leaves, fresh coriander, Chinese parsley, or cilantro (US, commercially in Canada, and Spanish-speaking countries). The fresh leaves are an ingredient in many foods, such as chutneys and salads, salsa, guacamole, and as a widely used garnish for soup, fish, and meat. As heat diminishes their flavour, coriander leaves are often used raw or added to the dish immediately before serving. In Indian and Central Asian recipes, coriander leaves are used in large amounts and cooked until the flavour diminishes. The leaves spoil quickly when removed from the plant and lose their aroma when dried or frozen.
The taste of the leaves differs from that of the seeds. The seeds exhibit citrus overtones. The dominant flavorants in the leaves are the aldehydes 2-decenal and 2-dodecenal. The main flavorant in the seeds is (+)-linalool.
Seeds
The dry fruits are coriander seeds. The word "coriander" in food preparation may refer solely to these seeds (as a spice), rather than the plant. The seeds have a lemony citrus flavour when crushed due to the terpenes linalool (which comprises about two thirds of its volatile components) and pinene. It is described as warm, nutty, spicy, and orange-flavoured.
The variety C. sativum var. sativum has a fruit diameter of , while var. microcarpum fruits have a diameter of , and var. indicum has elongated fruits. Large-fruited types are grown mainly by tropical and subtropical countries, such as Morocco, India, and Australia, and contain a low volatile oil content (0.1–0.4%). They are used for grinding and blending purposes in the spice trade. Types with smaller fruit are produced in temperate regions and usually have a volatile oil content of around 0.4–1.8%, so they are highly valued as a raw material for the preparation of essential oil.
Coriander is commonly found both as whole dried seeds and in ground form. Roasting or heating the seeds in a dry pan heightens the flavour, aroma, and pungency. Ground coriander seed loses flavour quickly in storage and is best ground fresh. Coriander seed is a spice in garam masala, and Indian curries, which often employ the ground fruits in generous amounts together with cumin, acting as a thickener in a mixture called dhania jeera. Roasted coriander seeds, called dhania dal, are eaten as a snack.
Outside of Asia, coriander seed is used widely for pickling vegetables. In Germany and South Africa (see boerewors), the seeds are used while making sausages. In Russia and Central Europe, coriander seed is an occasional ingredient in rye bread (e.g. Borodinsky bread) as an alternative to caraway. The Zuni people of North America have adapted it into their cuisine, mixing the powdered seeds ground with chilli, using it as a condiment with meat, and eating leaves as a salad.
Coriander seeds are used in brewing certain styles of beer, particularly some Belgian wheat beers. The coriander seeds are used with orange peel to add a citrus character.
Coriander seeds are one of the key botanicals used to flavour gin.
One preliminary study showed coriander essential oil to inhibit Gram-positive and Gram-negative bacteria, including Staphylococcus aureus, Enterococcus faecalis, Pseudomonas aeruginosa, and Escherichia coli.
Coriander is listed as one of the original ingredients in the secret formula for Coca-Cola.
Roots
Coriander roots have a deeper, more intense flavour than the leaves and are used in a variety of Asian cuisines, particularly in Thai dishes such as soups or curry pastes.
In culture
Coriander was mentioned by Hippocrates (around 400 BCE), as well as Dioscorides (65 CE).
| Biology and health sciences | Herbs and spices | Plants |
341682 | https://en.wikipedia.org/wiki/Square%20root%20of%202 | Square root of 2 | The square root of 2 (approximately 1.4142) is the positive real number that, when multiplied by itself or squared, equals the number 2. It may be written in mathematics as or . It is an algebraic number, and therefore not a transcendental number. Technically, it should be called the principal square root of 2, to distinguish it from the negative number with the same property.
Geometrically, the square root of 2 is the length of a diagonal across a square with sides of one unit of length; this follows from the Pythagorean theorem. It was probably the first number known to be irrational. The fraction (≈ 1.4142857) is sometimes used as a good rational approximation with a reasonably small denominator.
Sequence in the On-Line Encyclopedia of Integer Sequences consists of the digits in the decimal expansion of the square root of 2, here truncated to 65 decimal places:
History
The Babylonian clay tablet YBC 7289 (–1600 BC) gives an approximation of in four sexagesimal figures, , which is accurate to about six decimal digits, and is the closest possible three-place sexagesimal representation of , representing a margin of error of only –0.000042%:
Another early approximation is given in ancient Indian mathematical texts, the Sulbasutras (–200 BC), as follows: Increase the length [of the side] by its third and this third by its own fourth less the thirty-fourth part of that fourth. That is,
This approximation, diverging from the actual value of by approximately +0.07%, is the seventh in a sequence of increasingly accurate approximations based on the sequence of Pell numbers, which can be derived from the continued fraction expansion of . Despite having a smaller denominator, it is only slightly less accurate than the Babylonian approximation.
Pythagoreans discovered that the diagonal of a square is incommensurable with its side, or in modern language, that the square root of two is irrational. Little is known with certainty about the time or circumstances of this discovery, but the name of Hippasus of Metapontum is often mentioned. For a while, the Pythagoreans treated as an official secret the discovery that the square root of two is irrational, and, according to legend, Hippasus was murdered for divulging it, though this has little to any substantial evidence in traditional historian practice. The square root of two is occasionally called Pythagoras's number or Pythagoras's constant.
Ancient Roman architecture
In ancient Roman architecture, Vitruvius describes the use of the square root of 2 progression or ad quadratum technique. It consists basically in a geometric, rather than arithmetic, method to double a square, in which the diagonal of the original square is equal to the side of the resulting square. Vitruvius attributes the idea to Plato. The system was employed to build pavements by creating a square tangent to the corners of the original square at 45 degrees of it. The proportion was also used to design atria by giving them a length equal to a diagonal taken from a square, whose sides are equivalent to the intended atrium's width.
Decimal value
Computation algorithms
There are many algorithms for approximating as a ratio of integers or as a decimal. The most common algorithm for this, which is used as a basis in many computers and calculators, is the Babylonian method for computing square roots, an example of Newton's method for computing roots of arbitrary functions. It goes as follows:
First, pick a guess, ; the value of the guess affects only how many iterations are required to reach an approximation of a certain accuracy. Then, using that guess, iterate through the following recursive computation:
Each iteration improves the approximation, roughly doubling the number of correct digits. Starting with , the subsequent iterations yield:
Rational approximations
A simple rational approximation (≈ 1.4142857) is sometimes used. Despite having a denominator of only 70, it differs from the correct value by less than (approx. ).
The next two better rational approximations are (≈ 1.4141414...) with a marginally smaller error (approx. ), and (≈ 1.4142012) with an error of approx .
The rational approximation of the square root of two derived from four iterations of the Babylonian method after starting with () is too large by about ; its square is ≈ .
Records in computation
In 1997, the value of was calculated to 137,438,953,444 decimal places by Yasumasa Kanada's team. In February 2006, the record for the calculation of was eclipsed with the use of a home computer. Shigeru Kondo calculated one trillion decimal places in 2010. Other mathematical constants whose decimal expansions have been calculated to similarly high precision include , , and the golden ratio. Such computations provide empirical evidence of whether these numbers are normal.
This is a table of recent records in calculating the digits of .
Proofs of irrationality
Proof by infinite descent
One proof of the number's irrationality is the following proof by infinite descent. It is also a proof of a negation by refutation: it proves the statement " is not rational" by assuming that it is rational and then deriving a falsehood.
Assume that is a rational number, meaning that there exists a pair of integers whose ratio is exactly .
If the two integers have a common factor, it can be eliminated using the Euclidean algorithm.
Then can be written as an irreducible fraction such that and are coprime integers (having no common factor) which additionally means that at least one of or must be odd.
It follows that and . ( ) ( are integers)
Therefore, is even because it is equal to . ( is necessarily even because it is 2 times another whole number.)
It follows that must be even (as squares of odd integers are never even).
Because is even, there exists an integer that fulfills .
Substituting from step 7 for in the second equation of step 4: , which is equivalent to .
Because is divisible by two and therefore even, and because , it follows that is also even which means that is even.
By steps 5 and 8, and are both even, which contradicts step 3 (that is irreducible).
Since we have derived a falsehood, the assumption (1) that is a rational number must be false. This means that is not a rational number; that is to say, is irrational.
This proof was hinted at by Aristotle, in his Analytica Priora, §I.23. It appeared first as a full proof in Euclid's Elements, as proposition 117 of Book X. However, since the early 19th century, historians have agreed that this proof is an interpolation and not attributable to Euclid.
Proof using reciprocals
Assume by way of contradiction that were rational. Then we may write as an irreducible fraction in lowest terms, with coprime positive integers . Since , it follows that can be expressed as the irreducible fraction . However, since and differ by an integer, it follows that the denominators of their irreducible fraction representations must be the same, i.e. . This gives the desired contradiction.
Proof by unique factorization
As with the proof by infinite descent, we obtain . Being the same quantity, each side has the same prime factorization by the fundamental theorem of arithmetic, and in particular, would have to have the factor 2 occur the same number of times. However, the factor 2 appears an odd number of times on the right, but an even number of times on the left—a contradiction.
Application of the rational root theorem
The irrationality of also follows from the rational root theorem, which states that a rational root of a polynomial, if it exists, must be the quotient of a factor of the constant term and a factor of the leading coefficient. In the case of , the only possible rational roots are and . As is not equal to or , it follows that is irrational. This application also invokes the integer root theorem, a stronger version of the rational root theorem for the case when is a monic polynomial with integer coefficients; for such a polynomial, all roots are necessarily integers (which is not, as 2 is not a perfect square) or irrational.
The rational root theorem (or integer root theorem) may be used to show that any square root of any natural number that is not a perfect square is irrational. For other proofs that the square root of any non-square natural number is irrational, see Quadratic irrational number or Infinite descent.
Geometric proofs
A simple proof is attributed to Stanley Tennenbaum when he was a student in the early 1950s. Assume that , where and are coprime positive integers. Then and are the smallest positive integers for which . Now consider two squares with sides and , and place two copies of the smaller square inside the larger one as shown in Figure 1. The area of the square overlap region in the centre must equal the sum of the areas of the two uncovered squares. Hence there exist positive integers and such that . Since it can be seen geometrically that and , this contradicts the original assumption.
Tom M. Apostol made another geometric reductio ad absurdum argument showing that is irrational. It is also an example of proof by infinite descent. It makes use of classic compass and straightedge construction, proving the theorem by a method similar to that employed by ancient Greek geometers. It is essentially the same algebraic proof as in the previous paragraph, viewed geometrically in another way.
Let be a right isosceles triangle with hypotenuse length and legs as shown in Figure 2. By the Pythagorean theorem, . Suppose and are integers. Let be a ratio given in its lowest terms.
Draw the arcs and with centre . Join . It follows that , and and coincide. Therefore, the triangles and are congruent by SAS.
Because is a right angle and is half a right angle, is also a right isosceles triangle. Hence implies . By symmetry, , and is also a right isosceles triangle. It also follows that .
Hence, there is an even smaller right isosceles triangle, with hypotenuse length and legs . These values are integers even smaller than and and in the same ratio, contradicting the hypothesis that is in lowest terms. Therefore, and cannot be both integers; hence, is irrational.
Constructive proof
While the proofs by infinite descent are constructively valid when "irrational" is defined to mean "not rational", we can obtain a constructively stronger statement by using a positive definition of "irrational" as "quantifiably apart from every rational". Let and be positive integers such that (as satisfies these bounds). Now and cannot be equal, since the first has an odd number of factors 2 whereas the second has an even number of factors 2. Thus . Multiplying the absolute difference by in the numerator and denominator, we get
the latter inequality being true because it is assumed that , giving (otherwise the quantitative apartness can be trivially established). This gives a lower bound of for the difference , yielding a direct proof of irrationality in its constructively stronger form, not relying on the law of excluded middle. This proof constructively exhibits an explicit discrepancy between and any rational.
Proof by Pythagorean triples
This proof uses the following property of primitive Pythagorean triples:
If , , and are coprime positive integers such that , then is never even.
This lemma can be used to show that two identical perfect squares can never be added to produce another perfect square.
Suppose the contrary that is rational. Therefore,
where and
Squaring both sides,
Here, is a primitive Pythagorean triple, and from the lemma is never even. However, this contradicts the equation which implies that must be even.
Multiplicative inverse
The multiplicative inverse (reciprocal) of the square root of two is a widely used constant, with the decimal value:
It is often encountered in geometry and trigonometry because the unit vector, which makes a 45° angle with the axes in a plane, has the coordinates
Each coordinate satisfies
Properties
One interesting property of is
since
This is related to the property of silver ratios.
can also be expressed in terms of copies of the imaginary unit using only the square root and arithmetic operations, if the square root symbol is interpreted suitably for the complex numbers and :
is also the only real number other than 1 whose infinite tetrate (i.e., infinite exponential tower) is equal to its square. In other words: if for , and for , the limit of as will be called (if this limit exists) . Then is the only number for which . Or symbolically:
appears in Viète's formula for ,
which is related to the formula
Similar in appearance but with a finite number of terms, appears in various trigonometric constants:
It is not known whether is a normal number, which is a stronger property than irrationality, but statistical analyses of its binary expansion are consistent with the hypothesis that it is normal to base two.
Representations
Series and product
The identity , along with the infinite product representations for the sine and cosine, leads to products such as
and
or equivalently,
The number can also be expressed by taking the Taylor series of a trigonometric function. For example, the series for gives
The Taylor series of with and using the double factorial gives
The convergence of this series can be accelerated with an Euler transform, producing
It is not known whether can be represented with a BBP-type formula. BBP-type formulas are known for and , however.
The number can be represented by an infinite series of Egyptian fractions, with denominators defined by 2n th terms of a Fibonacci-like recurrence relation a(n) = 34a(n−1) − a(n−2), a(0) = 0, a(1) = 6.
Continued fraction
The square root of two has the following continued fraction representation:
The convergents formed by truncating this representation form a sequence of fractions that approximate the square root of two to increasing accuracy, and that are described by the Pell numbers (i.e., ). The first convergents are: and the convergent following is . The convergent differs from by almost exactly , which follows from:
Nested square
The following nested square expressions converge to
Applications
Paper size
In 1786, German physics professor Georg Christoph Lichtenberg found that any sheet of paper whose long edge is times longer than its short edge could be folded in half and aligned with its shorter side to produce a sheet with exactly the same proportions as the original. This ratio of lengths of the longer over the shorter side guarantees that cutting a sheet in half along a line results in the smaller sheets having the same (approximate) ratio as the original sheet. When Germany standardised paper sizes at the beginning of the 20th century, they used Lichtenberg's ratio to create the "A" series of paper sizes. Today, the (approximate) aspect ratio of paper sizes under ISO 216 (A4, A0, etc.) is 1:.
Proof:
Let shorter length and longer length of the sides of a sheet of paper, with
as required by ISO 216.
Let be the analogous ratio of the halved sheet, then
Physical sciences
There are some interesting properties involving the square root of 2 in the physical sciences:
The square root of two is the frequency ratio of a tritone interval in twelve-tone equal temperament music.
The square root of two forms the relationship of f-stops in photographic lenses, which in turn means that the ratio of areas between two successive apertures is 2.
The celestial latitude (declination) of the Sun during a planet's astronomical cross-quarter day points equals the tilt of the planet's axis divided by .
In the brain there are lattice cells, discovered in 2005 by a group led by May-Britt and Edvard Moser. "The grid cells were found in the cortical area located right next to the hippocampus [...] At one end of this cortical area the mesh size is small and at the other it is very large. However, the increase in mesh size is not left to chance, but increases by the squareroot of two from one area to the next."
| Mathematics | Basics | null |
341989 | https://en.wikipedia.org/wiki/Phasmatodea | Phasmatodea | The Phasmatodea (also known as Phasmida or Phasmatoptera) are an order of insects whose members are variously known as stick insects, stick bugs, walkingsticks, stick animals, or bug sticks. They are also occasionally referred to as Devil's darning needles, although this name is shared by both dragonflies and crane flies. They can be generally referred to as phasmatodeans, phasmids, or ghost insects, with phasmids in the family Phylliidae called leaf insects, leaf-bugs, walking leaves, or bug leaves. The group's name is derived from the Ancient Greek , meaning an apparition or phantom, referring to their resemblance to vegetation while in fact being animals. Their natural camouflage makes them difficult for predators to detect; still, many species have one of several secondary lines of defense in the form of startle displays, spines or toxic secretions. Stick insects from the genera Phryganistria, Ctenomorpha, and Phobaeticus include the world's longest insects.
Members of the order are found on all continents except Antarctica, but they are most abundant in the tropics and subtropics. They are herbivorous, with many species living unobtrusively in the tree canopy. They have an incomplete metamorphosis life cycle with three stages: egg, nymph and adult. Many phasmids are parthenogenic, and do not require fertilized eggs for female offspring to be produced. In hotter climates, they may breed all year round; in more temperate regions, the females lay eggs in the autumn before dying, and the new generation hatches in the spring. Some species have wings and can disperse by flying, while others are more restricted.
Description
Phasmids vary greatly in size, with females typically growing larger than males of the same species. Males of the smallest species, such as Timema cristinae, reach about long, while females of the longest, an undescribed species informally known as Phryganistria "chinensis", can be up to in total length, including outstretched legs. This makes it the world's longest insect. The heaviest species of phasmid is likely to be Heteropteryx dilatata, the females of which may weigh as much as .
Some phasmids have cylindrical stick-like shapes, while others have flattened, leaflike shapes. Many species are wingless, or have reduced wings. The thorax is long in the winged species, since it houses the flight muscles, and is typically much shorter in the wingless forms. Where present, the first pair of wings is narrow and cornified (hardened), while the hind wings are broad, with straight veins along their length and multiple cross-veins. Their wing venation is unique among insects. The body is often further modified to resemble vegetation, with ridges resembling leaf veins, bark-like tubercles, and other forms of camouflage. A few species, such as Carausius morosus, are even able to change their pigmentation to match their surroundings. The mouthparts project out from the head. Chewing mandibles are uniform across species. The legs are typically long and slender, and some species are capable of limb autotomy (appendage shedding). Phasmids have long, slender antennae, as long as or longer than the rest of the body in some species.
All phasmids possess compound eyes, but ocelli (light-sensitive organs) are only known from the five groups Lanceocercata, Necrosciinae, Pseudophasmatidae, Palophidae and Phylliidae. Of these only the first three groups have females with ocelli, which like the wings seems to have re-evolved from ancestors that had lost them. Phasmids have an impressive visual system that allows them to perceive significant detail even in dim conditions, which suits their typically nocturnal lifestyle. They are born equipped with tiny compound eyes with a limited number of facets. As phasmids grow through successive molts, the number of facets in each eye is increased along with the number of photoreceptor cells. The sensitivity of the adult eye is at least tenfold that of the nymph in its first instar (developmental stage). As the eye grows more complex, the mechanisms to adapt to dark/light changes are also enhanced: eyes in dark conditions evidence fewer screening pigments, which would block light, than during the daytime, and changes in the width of the retinal layer to adapt to changes in available light are significantly more pronounced in adults. The larger size of the adult insects' eyes makes them more prone to radiation damage. This explains why fully grown individuals are mostly nocturnal. Lessened sensitivity to light in the newly emerged insects helps them to escape from the leaf litter wherein they are hatched and move upward into the more brightly illuminated foliage. Young stick insects are diurnal (daytime) feeders and move around freely, expanding their foraging range.
Stick insects have two types of pads on their legs: sticky "toe pads" and non-stick "heel pads" a little further up their legs. The heel pads are covered in microscopic hairs which create strong friction at low pressure, enabling them to grip without having to be peeled energetically from the surface at each step. The sticky toe pads are used to provide additional grip when climbing but are not used on a level surface.
Distribution
Phasmatodea can be found all over the world except for the Antarctic and Patagonia. They are most numerous in the tropics and subtropics. The greatest diversity is found in Southeast Asia and South America, followed by Australia, Central America, and the southern United States. Over 300 species are known from the island of Borneo, making it the richest place in the world for Phasmatodea.
Anti-predator adaptations
Phasmatodea species exhibit mechanisms for defense from predators that prevent an attack from happening in the first place (primary defense), and defenses that are deployed after an attack has been initiated (secondary defense).
The defense mechanism most readily identifiable with Phasmatodea is camouflage, in the form of a plant mimicry. Most phasmids are known for effectively replicating the forms of sticks and leaves, and the bodies of some species (such as Pseudodiacantha macklotti and Bactrododema centaurum) are covered in mossy or lichenous outgrowths that supplement their disguise. Remaining absolutely stationary enhances their inconspicuousness. Some species have the ability to change color as their surroundings shift (Bostra scabrinota, Timema californica). In a further behavioral adaptation to supplement crypsis, a number of species perform a rocking motion where the body is swayed from side to side; this is thought to mimic the movement of leaves or twigs swaying in the breeze. Another method by which stick insects avoid predation and resemble twigs is by entering a cataleptic state, where the insect adopts a rigid, motionless posture that can be maintained for a long period. The nocturnal feeding habits of adults also help Phasmatodea to remain concealed from predators.
In a seemingly different method of defense, many species of Phasmatodea seek to startle the encroaching predator by flashing bright colors that are normally hidden, and making a loud noise. When disturbed on a branch or foliage, some species, while dropping to the undergrowth to escape, will open their wings momentarily during free fall to display bright colors that disappear when the insect lands. Others will maintain their display for up to 20 minutes, hoping to frighten the predator and convey the appearance of a larger size. Some, such as Pterinoxylus spinulosus, accompany the visual display with the noise made by rubbing together parts of the wings.
Some species, such as the young nymphs of Extatosoma tiaratum, have been observed to curl the abdomen upwards over the body and head to resemble ants or scorpions in an act of mimicry, another defense mechanism by which the insects avoid becoming prey. The eggs of some species such as Diapheromera femorata have fleshy projections resembling elaiosomes (fleshy structures sometimes attached to seeds) that attract ants. When the egg has been carried to the colony, the adult ant feeds the elaiosome to a larva while the phasmid egg is left to develop in the recesses of the nest in a protected environment.
When threatened, some phasmids that are equipped with femoral spines on the metathoracic legs (Oncotophasma martini, Eurycantha calcarata, Eurycantha horrida, Diapheromera veliei, Diapheromera covilleae, Heteropteryx dilatata) respond by curling the abdomen upward and repeatedly swinging the legs together, grasping at the threat. If the menace is caught, the spines can, in humans, draw blood and inflict considerable pain.
Some species are equipped with a pair of glands at the anterior (front) edge of the prothorax that enables the insect to release defensive secretions, including chemical compounds of varying effect: some produce distinct odors, and others can cause a stinging, burning sensation in the eyes and mouth of a predator. The spray often contains pungent-smelling volatile metabolites, previously thought to be concentrated in the insect from its plant food sources. However, it now seems more likely that the insect manufactures its own defensive chemicals. Additionally, the chemistry of the defense spray from at least one species, Anisomorpha buprestoides, has been shown to vary based on the insect's life stage or the particular population it is part of. This chemical spray variation also corresponds with regionally specific color forms in populations in Florida, with the different variants having distinct behaviors. The spray from one species, Megacrania nigrosulfurea, is used as a treatment for skin infections by a tribe in Papua New Guinea because of its antibacterial constituents. Some species employ a shorter-range defensive secretion, where individuals bleed reflexively through the joints of their legs and the seams of the exoskeleton when bothered, allowing the blood (hemolymph), which contains distasteful compounds, to discourage predators. Another ploy is to regurgitate their stomach contents when harassed, repelling potential predators.
Life cycle
The life cycle of the stick insect begins when the female deposits her eggs through one of these methods of oviposition: she will either flick her egg to the ground by a movement of the ovipositor or her entire abdomen, carefully place the eggs in the axils of the host plant, bury them in small pits in the soil, or stick the eggs to a substrate, usually a stem or leaf of the food plant. A single female lays from 100 to 1,200 eggs after mating, depending on the species.
Many species of phasmids are parthenogenic, meaning the females lay eggs without needing to mate with males to produce offspring. Eggs from virgin mothers are entirely female and hatch into nymphs that are exact copies of their mothers. Stick insect species that are the product of hybridisation are usually obligate parthenogens, but non-hybrids are facultative parthenogens, meaning they retain the ability to mate and their sexual behavior depends on the presence and abundance of males.
Phasmatodea eggs resemble seeds in shape and size and have hard shells. They have a lid-like structure called an operculum at the anterior pole, from which the nymph emerges during hatching. The eggs vary in the length of time before they hatch which varies from 13 to more than 70 days, with the average around 20 to 30 days. Some species, particularly those from temperate regions, undergo diapause, where development is delayed during the winter months. Diapause is initiated by the effect of short day lengths on the egg-laying adults or can be genetically determined. Diapause is broken by exposure to the cold of winter, causing the eggs to hatch during the following spring. Among species of economic importance such as Diapheromera femorata, diapause results in the development of two-year cycles of outbreaks.
Many species' eggs bear a fatty, knoblike capitulum that caps the operculum. This structure attracts ants because of its resemblance to the elaiosome of some plant seeds that are sought-after food sources for ant larvae, and usually contribute to ensuring seed dispersal by ants, a form of ant-plant mutualism called myrmecochory. The ants take the egg into their nest underground and can remove the capitulum to feed to their larvae without harming the phasmid embryo. There, the egg hatches and the young nymph, which initially resembles an ant (another instance of mimicry among Phasmatodea), eventually emerges from the nest and climbs the nearest tree to safety in the foliage. The eggs of stick insects have a coating of calcium oxalate which makes them survive unscathed in the digestive tract of birds. It has been suggested that birds may have a role in the dispersal of parthenogenetic stick insect species, especially to islands.
The Phasmatodea life cycle is hemimetabolous, proceeding through a series of several nymphal instars. Once emerged, a nymph will eat its cast skin. Adulthood is reached for most species after several months and many molts. The lifespan of Phasmatodea varies by species, but ranges from a few months to up to three years.
Ecology
Phasmids are herbivorous, feeding mostly on the leaves of trees and shrubs, and a conspicuous component of many Neotropical systems. Phasmatodea has been postulated as dominant light-gap herbivores there. Their role in the forest ecosystem is considered important by many scientists, who stress the significance of light gaps in maintaining succession and resilience in climax forests. The presence of phasmids lowers the net production of early successional plants by consuming them and then enriches the soil by defecation. This enables the late succession plants to become established and encourages the recycling of the tropical forest.
Phasmatodea are recognized as injurious to forest and shade trees by defoliation. Didymuria violescens, Podacanthus wilkinsoni and Ctenomorphodes tessulatus in Australia, Diapheromera femorata in North America and Graeffea crouani in coconut plantations in the South Pacific all occur in outbreaks of economic importance. Indeed, in the American South, as well as in Michigan and Wisconsin, the walking stick is a significant problem in parks and recreation sites, where it consumes the foliage of oaks and other hardwoods. Severe outbreaks of the walking stick, Diapheromera femorata, have occurred in the Ouachita Mountains of Arkansas and Oklahoma. The insects eat the entire leaf blade. In the event of heavy outbreaks, entire stands of trees can be completely denuded. Continuous defoliation over several years often results in the death of the tree. Because these species cannot fly, infestations are typically contained to a radius of a few hundred yards. Nevertheless, the damage incurred to parks in the region is often costly. Control efforts in the case of infestations have typically involved chemical pesticides; ground fires are effective at killing eggs but have obvious disadvantages. In New South Wales, research has investigated the feasibility of controlling stick insects using natural enemies such as parasitic wasps (Myrmecomimesis spp.).
Taxonomy
The classification of the Phasmatodea is complex and the relationships between its members are poorly understood. Furthermore, there is much confusion over the ordinal name. Phasmida is preferred by many authors, though it is incorrectly formed; Phasmatodea is correctly formed, and is widely accepted. However, Brock and Marshall argue:
The order Phasmatodea is sometimes considered to be related to other orders, including the Blattodea, Mantodea, Grylloblattodea, Mantophasmatodea and Dermaptera, but the affiliations are uncertain and the grouping (sometimes referred to as "Orthopteroidea") may be paraphyletic (not have a common ancestor) and hence invalid in the traditional circumscription (set of attributes that all members have). Phasmatodea, once considered a suborder of Orthoptera, is now treated as an order of its own. Anatomical features separate them as a monophyletic (descended from a common ancestor) group from the Orthoptera. One is the instance among all species of Phasmatodea of a pair of exocrine glands inside the prothorax used for defense. Another is the presence of a specially formed sclerite (hardened plate), called a vomer, which allows the male to clasp the female during mating.
The order is divided into two, or sometimes three, suborders. The traditional division is into the suborder groups Anareolatae and Areolatae, which are distinguished according to whether the insect has sunken areola, or circular areas, on the underside of the apices of the middle and hind tibiae (Areolate) or not (Anareolate). However the phylogenetic (evolutionary) relationships between the different groups is poorly resolved. The monophyly of Anareolatae has been questioned and the morphology of the eggs may be a better basis for classification. An alternative is to divide the Phasmatodea into three suborders Agathemerodea (1 genus and 8 species), Timematodea (1 genus and 21 species) and Euphasmatodea (or Verophasmatodea) for the remaining taxa. This division is, however, not fully supported by the molecular studies, which recover Agathemerodea as nested within Euphasmatodea rather than being the sister group of the latter group. Recent taxonomic treatments recognise two suborders, with Agathemeridae placed in Pseudophasmatoidea within Euphasmatodea and Agathemerodea treated as nomen dubium
While suggestions have been made that various insects extending back to the Permian epoch represent stem-group phasmatodeans, the earliest unambiguous members of the group are the Susumanioidea, which first appeared during the Middle Jurassic, and usually have two large pairs of wings. Modern phasmatodeans first appeared during the Early Cretaceous, with the cuirrently oldest known being Araripephasma from the Early Cretaceous (Aptian) Crato Formation of Brazil, around 113 million years old, which can be confidently assigned to the Euphasmatodea.
The earliest leaf insect (Phylliinae) fossil is Eophyllium messelensis from the 47-million-year-old Eocene of Messel, Germany. In size and cryptic (leaflike) body form, it closely resembles extant species, suggesting that the behavior of the group has changed little since that time.
Over 3,500 species have been described, with many more yet to be described both in museum collections and in the wild.
Select species
One Australian species, the Lord Howe Island stick insect, is now listed as critically endangered. It was believed extinct until its rediscovery on the rock known as Ball's Pyramid. An effort is underway in Australia to rear this species in captivity.
The best known of the stick insects is the Indian or laboratory stick insect (Carausius morosus). This insect grows to roughly 10 cm (4 in) and reproduces parthenogenically, and although males have been recorded, they are rare.
Fossils of the extinct genus and species Eoprephasma hichensi have been recovered from Ypresian age sediments in the U.S. state of Washington and British Columbia, Canada. The species is one of the youngest members of the stem phasmatodean group Susumanioidea.
European species
In Europe there are 17 species of stick insects described, belonging to the genera Bacillus Clonopsis, Leptynia and Pijnackeria. There are also a few other species that live in Europe but are introduced, as for example with a couple of species of Acanthoxyla, which are native to New Zealand but are present in southern England.
In the Iberian Peninsula there are currently described 13 species and several subspecies. Their life cycle is annual, living only during the hottest months (especially genera Leptynia and Pijnackeria), which usually means late spring to early autumn.
Behavior
Stick insects, like praying mantises, show rocking behavior in which the insect makes rhythmic, repetitive, side-to-side movements. The common interpretation of this behavior's function is it enhances crypsis by mimicking vegetation moving in the wind. These movements may also be important in allowing the insects to discriminate objects from the background by relative motion. Rocking movements by these generally sedentary insects may replace flying or running as a source of relative motion to help them discern objects in the foreground.
Mating behavior in Phasmatodea is impressive because of the extraordinarily long duration of some pairings. A record among insects, the stick insect Necroscia sparaxes, found in India, is sometimes coupled for 79 days at a time. It is not uncommon for this species to assume the mating posture for days or weeks on end, and among some species (Diapheromera veliei and D. covilleae), pairing can last three to 136 hours in captivity.
Overt displays of aggression between males over mates suggests that extended pairing may have evolved to guard females from sperm competition. Fighting between competing males has been observed in the species D. veiliei and D. covilleae. During these encounters, the approach of a challenger causes the existing mate to manipulate the female's abdomen, which he has clasped by means of the clasping organ, or vomer, down upon itself to block the site of attachment. Occasionally, the consort will strike out at the competitor with the mid femora, which are equipped with an enlarged and hooked spine in both sexes that can draw the blood of the opponent when they are flexed against the body to puncture the integument. Usually, a strong hold on the female's abdomen and blows to the intruder are enough to deter the unwanted competition, but occasionally the competitor has been observed to employ a sneaky tactic to inseminate the female. While the first mate is engaged in feeding and is forced to vacate the dorsal position, the intruder can clasp the female's abdomen and insert his genitalia. If he is discovered, the males will enter into combat wherein they lean backward, both clasped to the female's abdomen, and freely suspended, engage in rapid, sweeping blows with their forelegs in a manner similar to boxing. Usually, when the intruder gains attachment to the female's abdomen, these conflicts result in the displacement of the original mate.
Lengthy pairings have also been described in terms of a defensive alliance. When cleaved together, the pair is more unwieldy for predators to handle. Also, the chemical defenses (secretions, reflex bleeding, regurgitation) of the individual stick insect are enhanced when two are paired. Females survive attacks by predators significantly better when pairing, largely because the dorsal position of the male functions well as a shield. This could indicate that manipulation by females is taking place: if females accept ejaculate at a slow rate, for instance, the males are forced to remain in copulo for longer and the female's chances of survival are enhanced. Also, evolution could have simply favored males that remained attached to their females longer, since females are often less abundant than males and represent a valuable prize, so for the lucky male, even the sacrifice of his own life to preserve his offspring with the female may be worthwhile. Sexual dimorphism in the species, where females are usually significantly larger than the males, may have evolved due to the fitness advantage accrued to males that can remain attached to the female, thereby blocking competitors, without severely impeding her movement.
Certain Phasmatodea, such as Anisomorpha buprestoides, sometimes form aggregations. These insects have been observed to congregate during the day in a concealed location, going their separate ways at nightfall to forage, and returning to their refuge before dawn. Such behavior has been little studied, and how the insects find their way back is unknown.
In human culture
Stick insects are often kept in captivity: almost 300 species have been reared in laboratories or as pets. The most commonly kept is the Indian (or laboratory) stick insect, Carausius morosus, which eats vegetables such as lettuce. Droppings of the stick insect Eurycnema versirubra (Serville, 1838) [=Eurycnema versifasciata] fed with specific plants are made into a medicinal tea by Malaysian Chinese to treat ailments.
The botanical illustrator Marianne North (1830–1890) painted leaf and stick insects that she saw on her travels in the 1870s.
Tribesmen in Sarawak eat phasmids and their eggs.
Some indigenous people of the D'Entrecasteaux Islands have traditionally made fishhooks from the legs of certain phasmids.
Research has been conducted to analyze the stick insect method of walking and apply this to the engineering of six-legged walking robots. Instead of one centralized control system, it seems each leg of a phasmid operates independently.
In Australia and Hawaii many kinds of stick insects are kept as exotic pets including the strong, goliath, spiny and children's. The custom of keeping stick insects as pets was probably brought to Australia by either Chinese, Japanese or Vietnamese immigrants during World War II, the Korean War or the Vietnam War.
Stick insects have been kept as pets since the time of the Han dynasty. They were kept inside birdcages and people in the Far East believe they bring good luck and fortune, just like crickets.
The video game Disco Elysium includes a storyline centered around a giant stick insect and cryptid called the insulindian phasmid.
A clip of a stick insect swaying back and forth, in a manner akin to dancing, became an Internet meme in 2020 as a bait-and-switch.
| Biology and health sciences | Insects and other hexapods | null |
342078 | https://en.wikipedia.org/wiki/Gimbal | Gimbal | A gimbal is a pivoted support that permits rotation of an object about an axis. A set of three gimbals, one mounted on the other with orthogonal pivot axes, may be used to allow an object mounted on the innermost gimbal to remain independent of the rotation of its support (e.g. vertical in the first animation). For example, on a ship, the gyroscopes, shipboard compasses, stoves, and even drink holders typically use gimbals to keep them upright with respect to the horizon despite the ship's pitching and rolling.
The gimbal suspension used for mounting compasses and the like is sometimes called a Cardan suspension after Italian mathematician and physicist Gerolamo Cardano (1501–1576) who described it in detail. However, Cardano did not invent the gimbal, nor did he claim to. The device has been known since antiquity, first described in the 3rd c. BC by Philo of Byzantium, although some modern authors support the view that it may not have a single identifiable inventor.
History
The gimbal was first described by the Greek inventor Philo of Byzantium (280–220 BC). Philo described an eight-sided ink pot with an opening on each side, which can be turned so that while any face is on top, a pen can be dipped and inked — yet the ink never runs out through the holes of the other sides. This was done by the suspension of the inkwell at the center, which was mounted on a series of concentric metal rings so that it remained stationary no matter which way the pot is turned.
In Ancient China, the Han dynasty (202 BC – 220 AD) inventor and mechanical engineer Ding Huan created a gimbal incense burner around 180 AD. There is a hint in the writing of the earlier Sima Xiangru (179–117 BC) that the gimbal existed in China since the 2nd century BC. There is mention during the Liang dynasty (502–557) that gimbals were used for hinges of doors and windows, while an artisan once presented a portable warming stove to Empress Wu Zetian (r. 690–705) which employed gimbals. Extant specimens of Chinese gimbals used for incense burners date to the early Tang dynasty (618–907), and were part of the silver-smithing tradition in China.
The authenticity of Philo's description of a cardan suspension has been doubted by some authors on the ground that the part of Philo's Pneumatica which describes the use of the gimbal survived only in an Arabic translation of the early 9th century. Thus, as late as 1965, the sinologist Joseph Needham suspected Arab interpolation. However, Carra de Vaux, author of the French translation which still provides the basis for modern scholars, regards the Pneumatics as essentially genuine. The historian of technology George Sarton (1959) also asserts that it is safe to assume the Arabic version is a faithful copying of Philo's original, and credits Philon explicitly with the invention. So does his colleague Michael Lewis (2001). In fact, research by the latter scholar (1997) demonstrates that the Arab copy contains sequences of Greek letters which fell out of use after the 1st century, thereby strengthening the case that it is a faithful copy of the Hellenistic original, a view recently also shared by the classicist Andrew Wilson (2002).
The ancient Roman author Athenaeus Mechanicus, writing during the reign of Augustus (30 BC–14 AD), described the military use of a gimbal-like mechanism, calling it "little ape" (pithêkion). When preparing to attack coastal towns from the sea-side, military engineers used to yoke merchant-ships together to take the siege machines up to the walls. But to prevent the shipborne machinery from rolling around the deck in heavy seas, Athenaeus advises that "you must fix the pithêkion on the platform attached to the merchant-ships in the middle, so that the machine stays upright in any angle".
After antiquity, gimbals remained widely known in the Near East. In the Latin West, reference to the device appeared again in the 9th century recipe book called the Little Key of Painting''' (mappae clavicula). The French inventor Villard de Honnecourt depicts a set of gimbals in his sketchbook (see right). In the early modern period, dry compasses were suspended in gimbals.
Applications
Inertial navigation
In inertial navigation, as applied to ships and submarines, a minimum of three gimbals are needed to allow an inertial navigation system (stable table) to remain fixed in inertial space, compensating for changes in the ship's yaw, pitch, and roll. In this application, the inertial measurement unit (IMU) is equipped with three orthogonally mounted gyros to sense rotation about all axes in three-dimensional space. The gyro outputs are kept to a null through drive motors on each gimbal axis, to maintain the orientation of the IMU. To accomplish this, the gyro error signals are passed through "resolvers" mounted on the three gimbals, roll, pitch and yaw. These resolvers perform an automatic matrix transformation according to each gimbal angle, so that the required torques are delivered to the appropriate gimbal axis. The yaw torques must be resolved by roll and pitch transformations. The gimbal angle is never measured.
Similar sensing platforms are used on aircraft.
In inertial navigation systems, gimbal lock may occur when vehicle rotation causes two of the three gimbal rings to align with their pivot axes in a single plane. When this occurs, it is no longer possible to maintain the sensing platform's orientation.
Rocket engines
In spacecraft propulsion, rocket engines are generally mounted on a pair of gimbals to allow a single engine to vector thrust about both the pitch and yaw axes; or sometimes just one axis is provided per engine. To control roll, twin engines with differential pitch or yaw control signals are used to provide torque about the vehicle's roll axis.
Photography and imaging
Gimbals are also used to mount everything from small camera lenses to large photographic telescopes.
In portable photography equipment, single-axis gimbal heads are used in order to allow a balanced movement for camera and lenses. This proves useful in wildlife photography as well as in any other case where very long and heavy telephoto lenses are adopted: a gimbal head rotates a lens around its center of gravity, thus allowing for easy and smooth manipulation while tracking moving subjects.
Very large gimbal mounts in the form 2 or 3 axis altitude-altitude mounts are used in satellite photography for tracking purposes.
Gyrostabilized gimbals which house multiple sensors are also used for airborne surveillance applications including airborne law enforcement, pipe and power line inspection, mapping, and ISR (intelligence, surveillance, and reconnaissance). Sensors include thermal imaging, daylight, low light cameras as well as laser range finder, and illuminators.
Gimbal systems are also used in scientific optics equipment. For example, they are used to rotate a material sample along an axis to study their angular dependence of optical properties.
Film and video
Handheld 3-axis gimbals are used in stabilization systems designed to give the camera operator the independence of handheld shooting without camera vibration or shake. There are two versions of such stabilization systems: mechanical and motorized.
Mechanical gimbals have the sled, which includes the top stage where the camera is attached, the post'' which in most models can be extended, with the monitor and batteries at the bottom to counterbalance the camera weight. This is how the Steadicam stays upright, by simply making the bottom slightly heavier than the top, pivoting at the gimbal. This leaves the center of gravity of the whole rig, however heavy it may be, exactly at the operator's fingertip, allowing deft and finite control of the whole system with the lightest of touches on the gimbal.
Powered by three brushless motors, motorized gimbals have the ability to keep the camera level on all axes as the camera operator moves the camera. An inertial measurement unit (IMU) responds to movement and utilizes its three separate motors to stabilize the camera. With the guidance of algorithms, the stabilizer is able to notice the difference between deliberate movement such as pans and tracking shots from unwanted shake. This allows the camera to seem as if it is floating through the air, an effect achieved by a Steadicam in the past. Gimbals can be mounted to cars and other vehicles such as drones, where vibrations or other unexpected movements would make tripods or other camera mounts unacceptable. An example which is popular in the live TV broadcast industry, is the Newton 3-axis camera gimbal.
Marine chronometers
The rate of a mechanical marine chronometer is sensitive to its orientation. Because of this, chronometers were normally mounted on gimbals, in order to isolate them from the rocking motions of a ship at sea.
Gimbal lock
Gimbal lock is the loss of one degree of freedom in a three-dimensional, three-gimbal mechanism that occurs when the axes of two of the three gimbals are driven into a parallel configuration, "locking" the system into rotation in a degenerate two-dimensional space.
The word lock is misleading: no gimbal is restrained. All three gimbals can still rotate freely about their respective axes of suspension. Nevertheless, because of the parallel orientation of two of the gimbals' axes there is no gimbal available to accommodate rotation about one axis.
| Technology | Mechanisms | null |
342304 | https://en.wikipedia.org/wiki/Lumbar%20puncture | Lumbar puncture | Lumbar puncture (LP), also known as a spinal tap, is a medical procedure in which a needle is inserted into the spinal canal, most commonly to collect cerebrospinal fluid (CSF) for diagnostic testing. The main reason for a lumbar puncture is to help diagnose diseases of the central nervous system, including the brain and spine. Examples of these conditions include meningitis and subarachnoid hemorrhage. It may also be used therapeutically in some conditions. Increased intracranial pressure (pressure in the skull) is a contraindication, due to risk of brain matter being compressed and pushed toward the spine. Sometimes, lumbar puncture cannot be performed safely (for example due to a severe bleeding tendency). It is regarded as a safe procedure, but post-dural-puncture headache is a common side effect if a small atraumatic needle is not used.
The procedure is typically performed under local anesthesia using a sterile technique. A hypodermic needle is used to access the subarachnoid space and collect fluid. Fluid may be sent for biochemical, microbiological, and cytological analysis. Using ultrasound to landmark may increase success.
Lumbar puncture was first introduced in 1891 by the German physician Heinrich Quincke.
Medical uses
The reason for a lumbar puncture may be to make a diagnosis or to treat a disease, as outlined below.
Diagnosis
The chief diagnostic indications of lumbar puncture are for collection of cerebrospinal fluid (CSF). Analysis of CSF may exclude infectious, inflammatory, and neoplastic diseases affecting the central nervous system. The most common purpose is in suspected meningitis, since there is no other reliable tool with which meningitis, a life-threatening but highly treatable condition, can be excluded. A lumbar puncture can also be used to detect whether someone has Stage 1 or Stage 2 Trypanosoma brucei. Young infants commonly require lumbar puncture as a part of the routine workup for fever without a source. This is due to higher rates of meningitis than in older persons. Infants also do not reliably show classic symptoms of meningeal irritation (meningismus) like neck stiffness and headache the way adults do. In any age group, subarachnoid hemorrhage, hydrocephalus, benign intracranial hypertension, and many other diagnoses may be supported or excluded with this test. It may also be used to detect the presence of malignant cells in the CSF, as in carcinomatous meningitis or medulloblastoma. CSF containing less than 10 red blood cells (RBCs)/mm3 constitutes a "negative" tap in the context of a workup for subarachnoid hemorrhage, for example. Taps that are "positive" have an RBC count of 100/mm3 or more.
Treatment
Lumbar punctures may also be done to inject medications into the cerebrospinal fluid ("intrathecally"), particularly for spinal anesthesia or chemotherapy.
Serial lumbar punctures may be useful in temporary treatment of idiopathic intracranial hypertension (IIH). This disease is characterized by increased pressure of CSF which may cause headache and permanent loss of vision. While mainstays of treatment are medication, in some cases lumbar puncture performed multiple times may improve symptoms. It is not recommended as a staple of treatment due to discomfort and risk of the procedure, and the short duration of its efficacy.
Additionally, some people with normal pressure hydrocephalus (characterized by urinary incontinence, a changed ability to walk properly, and dementia) receive some relief of symptoms after removal of CSF.
Contraindications
Lumbar puncture should not be performed in the following situations:
Idiopathic (unidentified cause) increased intracranial pressure (ICP)
Rationale: lumbar puncture in the presence of raised ICP may cause uncal herniation
Exception: therapeutic use of lumbar puncture to reduce ICP, but only if obstruction (for example in the third ventricle of the brain) has been ruled out
Precaution
CT brain, especially in the following situations
Age >65
Reduced GCS
Recent history of seizure
Focal neurological signs
Abnormal respiratory pattern
Hypertension with bradycardia and deteriorating consciousness
Ophthalmoscopy for papilledema
Bleeding diathesis (relative)
Coagulopathy
Decreased platelet count (<50 x 109/L)
Infections
Skin infection at puncture site
Vertebral deformities (scoliosis or kyphosis), in hands of an inexperienced physician.
Adverse effects
Headache
Post-dural-puncture headache with nausea is the most common complication; it often responds to pain medications and infusion of fluids. It was long taught that this complication can be prevented by strict maintenance of a supine posture for two hours after the successful puncture; this has not been borne out in modern studies involving large numbers of people. Doing the procedure with the person on their side might decrease the risk. Intravenous caffeine injection is often quite effective in aborting these spinal headaches. A headache that is persistent despite a long period of bedrest and occurs only when sitting up may be indicative of a CSF leak from the lumbar puncture site. It can be treated by more bedrest, or by an epidural blood patch, where the person's own blood is injected back into the site of leakage to cause a clot to form and seal off the leak.
The risk of headache and need for analgesia and blood patch is much reduced if "atraumatic" needles are used. This does not affect the success rate of the procedure in other ways. Although the cost and difficulty are similar, adoption remains low, at only 16% .
The headaches may be caused by inadvertent puncture of the dura mater.
Other
Contact between the side of the lumbar puncture needle and a spinal nerve root can result in anomalous sensations (paresthesia) in a leg during the procedure; this is harmless and people can be warned about it in advance to minimize their anxiety if it should occur.
Serious complications of a properly performed lumbar puncture are extremely rare. They include spinal or epidural bleeding, adhesive arachnoiditis and trauma to the spinal cord or spinal nerve roots resulting in weakness or loss of sensation, or even paraplegia. The latter is exceedingly rare, since the level at which the spinal cord ends (normally the inferior border of L1, although it is slightly lower in infants) is several vertebral spaces above the proper location for a lumbar puncture (L3/L4). There are case reports of lumbar puncture resulting in perforation of abnormal dural arterio-venous malformations, resulting in catastrophic epidural hemorrhage; this is exceedingly rare.
The procedure is not recommended when epidural infection is present or suspected, when topical infections or dermatological conditions pose a risk of infection at the puncture site or in patients with severe psychosis or neurosis with back pain. Some authorities believe that withdrawal of fluid when initial pressures are abnormal could result in spinal cord compression or cerebral herniation; others believe that such events are merely coincidental in time, occurring independently as a result of the same pathology that the lumbar puncture was performed to diagnose. In any case, computed tomography of the brain is often performed prior to lumbar puncture if an intracranial mass is suspected.
CSF leaks can result from a lumbar puncture procedure.
Technique
Mechanism
The brain and spinal cord are enveloped by a layer of cerebrospinal fluid, 125–150 mL in total (in adults) which acts as a shock absorber and provides a medium for the transfer of nutrients and waste products. The majority is produced by the choroid plexus in the brain and circulates from there to other areas, before being reabsorbed into the circulation (predominantly by the arachnoid granulations).
The cerebrospinal fluid can be accessed most safely in the lumbar cistern. Below the first or second lumbar vertebrae (L1 or L2) the spinal cord terminates (conus medullaris). Nerves continue down the spine below this, but in a loose bundle of nerve fibers called the cauda equina. There is lower risk with inserting a needle into the spine at the level of the cauda equina because these loose fibers move out of the way of the needle without being damaged. The lumbar cistern extends into the sacrum up to the S2 vertebra.
Procedure
The person is usually placed on their side (left more commonly than right). The patient bends the neck so the chin is close to the chest, hunches the back, and brings knees toward the chest. This approximates a fetal position as much as possible. Patients may also sit on a stool and bend their head and shoulders forward. The area around the lower back is prepared using aseptic technique. Once the appropriate location is palpated, local anaesthetic is infiltrated under the skin and then injected along the intended path of the spinal needle. A spinal needle is inserted between the lumbar vertebrae L3/L4, L4/L5 or L5/S1 and pushed in until there is a "give" as it enters the lumbar cistern wherein the ligamentum flavum is housed. The needle is again pushed until there is a second 'give' that indicates the needle is now past the dura mater. The arachnoid membrane and the dura mater exist in flush contact with one another in the living person's spine due to fluid pressure from CSF in the subarachnoid space pushing the arachnoid membrane out towards the dura. Therefore, once the needle has pierced the dura mater it has also traversed the thinner arachnoid membrane. The needle is then in the subarachnoid space. The stylet from the spinal needle is then withdrawn and drops of cerebrospinal fluid are collected. The opening pressure of the cerebrospinal fluid may be taken during this collection by using a simple column manometer. The procedure is ended by withdrawing the needle while placing pressure on the puncture site. The spinal level is so selected to avoid spinal injuries. In the past, the patient would lie on their back for at least six hours and be monitored for signs of neurological problems. There is no scientific evidence that this provides any benefit. The technique described is almost identical to that used in spinal anesthesia, except that spinal anesthesia is more often done with the patient in a seated position.
The upright seated position is advantageous in that there is less distortion of spinal anatomy which allows for easier withdrawal of fluid. Some practitioners prefer it for lumbar puncture in obese patients, where lying on their side would cause a scoliosis and unreliable anatomical landmarks. However, opening pressures are notoriously unreliable when measured in the seated position. Therefore, patients will ideally lie on their side if practitioners need to measure opening pressure.
Reinsertion of the stylet may decrease the rate of post lumbar puncture headaches.
Although not available in all clinical settings, use of ultrasound is helpful for visualizing the interspinous space and assessing the depth of the spine from the skin. Use of ultrasound reduces the number of needle insertions and redirections, and results in higher rates of successful lumbar puncture. If the procedure is difficult, such as in people with spinal deformities such as scoliosis, it can also be performed under fluoroscopy (under continuous X-ray imaging).
Children
In children, a sitting flexed position was as successful as lying on the side with respect to obtaining non-traumatic CSF, CSF for culture, and cell count. There was a higher success rate in obtaining CSF in the first attempt in infants younger than 12 months in the sitting flexed position.
The spine of an infant at the time of birth differs from the adult spine. The conus medullaris (bottom of the spinal cord) terminates at the level of L1 in adults, but may range in term neonates (newly born babies) from L1–L3 levels. It is important to insert the spinal needle below the conus medullaris at the L3/L4 or L4/L5 interspinous levels. With growth of the spine, the conus typically reaches the adult level (L1) by 2 years of age.
The ligamentum flavum and dura mater are not as thick in infants and children as they are in adults. Therefore, it is difficult to assess when the needle passes through them into the subarachnoid space because the characteristic "pop" or "give" may be subtle or nonexistent in the pediatric lumbar puncture. To decrease the chances of inserting the spinal needle too far, some clinicians use the "Cincinnati" method. This method involves removing the stylet of the spinal needle once the needle has advanced through the dermis. After removal of the stylet, the needle is inserted until CSF starts to come out of the needle. Once all of the CSF is collected, the stylet is then reinserted before removal of the needle.
Newborn infants
Lumbar punctures are often used to diagnose or verify an infection in very young babies and can cause quite a bit of pain unless appropriate pain control is used (analgesia). Managing pain is important for infants undergoing this procedure. Approaches for pain control include topical pain medications (anaesthetics such as lidocaine). The most effective approach for pain control in infants who require a lumbar puncture is not clear.
Interpretation
Analysis of the cerebrospinal fluid generally includes a cell count and determination of the glucose and protein concentrations. The other analytical studies of cerebrospinal fluid are conducted according to the diagnostic suspicion.
Pressure determination
Increased CSF pressure can indicate congestive heart failure, cerebral edema, subarachnoid hemorrhage, hypo-osmolality resulting from hemodialysis, meningeal inflammation, purulent meningitis or tuberculous meningitis, hydrocephalus, or pseudotumor cerebri. In the setting of raised pressure (or normal pressure hydrocephalus, where the pressure is normal but there is excessive CSF), lumbar puncture may be therapeutic.
Decreased CSF pressure can indicate complete subarachnoid blockage, leakage of spinal fluid, severe dehydration, hyperosmolality, or circulatory collapse. Significant changes in pressure during the procedure can indicate tumors or spinal blockage resulting in a large pool of CSF, or hydrocephalus associated with large volumes of CSF.
Cell count
The presence of white blood cells in cerebrospinal fluid is called pleocytosis. A small number of monocytes can be normal; the presence of granulocytes is always an abnormal finding. A large number of granulocytes often heralds bacterial meningitis. White cells can also indicate reaction to repeated lumbar punctures, reactions to prior injections of medicines or dyes, central nervous system hemorrhage, leukemia, recent epileptic seizure, or a metastatic tumor. When peripheral blood contaminates the withdrawn CSF, a common procedural complication, white blood cells will be present along with erythrocytes, and their ratio will be the same as that in the peripheral blood.
The finding of erythrophagocytosis, where phagocytosed erythrocytes are observed, signifies haemorrhage into the CSF that preceded the lumbar puncture. Therefore, when erythrocytes are detected in the CSF sample, erythrophagocytosis suggests causes other than a traumatic tap, such as intracranial haemorrhage and haemorrhagic herpetic encephalitis. In which case, further investigations are warranted, including imaging and viral culture.
Microbiology
CSF can be sent to the microbiology lab for various types of smears and cultures to diagnose infections.
Gram staining may demonstrate gram positive bacteria in bacterial meningitis.
Microbiological culture is the gold standard for detecting bacterial meningitis. Bacteria, fungi, and viruses can all be cultured by using different techniques.
Polymerase chain reaction (PCR) has been a great advance in the diagnosis of some types of meningitis, such as meningitis from herpesvirus and enterovirus. It has high sensitivity and specificity for many infections of the CNS, is fast, and can be done with small volumes of CSF. Even though testing is expensive, cost analyses of PCR testing in neonatal patients demonstrated savings via reduced cost of hospitalization.
Numerous antibody-mediated tests for CSF are available in some countries: these include rapid tests for antigens of common bacterial pathogens, treponemal titers for the diagnosis of neurosyphilis and Lyme disease, Coccidioides antibody, and others.
The India ink test is still used for detection of meningitis caused by Cryptococcus neoformans, but the cryptococcal antigen (CrAg) test has a higher sensitivity.
Chemistry
Several substances found in cerebrospinal fluid are available for diagnostic measurement.
Glucose is present in the CSF; the level is usually about 60% that in the peripheral circulation. A fingerstick or venipuncture at the time of lumbar puncture may therefore be performed to assess peripheral glucose levels and determine a predicted CSF glucose value. Decreased glucose levels can indicate fungal, tuberculous or pyogenic infections; lymphomas; leukemia spreading to the meninges; meningoencephalitic mumps; or hypoglycemia. A glucose level of less than one third of blood glucose levels in association with low CSF lactate levels is typical in hereditary CSF glucose transporter deficiency also known as De Vivo disease.
Increased glucose levels in the fluid can indicate diabetes, although the 60% rule still applies.
Increased levels of glutamine are often involved with hepatic encephalopathies, Reye's syndrome, hepatic coma, cirrhosis, hypercapnia and depression.
Increased levels of lactate can occur the presence of cancer of the CNS, multiple sclerosis, heritable mitochondrial disease, low blood pressure, low serum phosphorus, respiratory alkalosis, idiopathic seizures, traumatic brain injury, cerebral ischemia, brain abscess, hydrocephalus, hypocapnia or bacterial meningitis.
The enzyme lactate dehydrogenase can be measured to help distinguish meningitides of bacterial origin, which are often associated with high levels of the enzyme, from those of viral origin in which the enzyme is low or absent.
Changes in total protein content of cerebrospinal fluid can result from pathologically increased permeability of the blood-cerebrospinal fluid barrier, obstructions of CSF circulation, meningitis, neurosyphilis, brain abscesses, subarachnoid hemorrhage, polio, collagen disease or Guillain–Barré syndrome, leakage of CSF, increases in intracranial pressure, or hyperthyroidism. Very high levels of protein may indicate tuberculous meningitis or spinal block.
IgG synthetic rate is calculated from measured IgG and total protein levels; it is elevated in immune disorders such as multiple sclerosis, transverse myelitis, and neuromyelitis optica of Devic. Oligoclonal bands may be detected in CSF but not in serum, suggesting intrathecal antibody production.
History
The first technique for accessing the dural space was described by the London physician Walter Essex Wynter. In 1889 he developed a crude cut down with cannulation in four patients with tuberculous meningitis. The main purpose was the treatment of raised intracranial pressure rather than for diagnosis. The technique for needle lumbar puncture was then introduced by the German physician Heinrich Quincke, who credits Wynter with the earlier discovery; he first reported his experiences at an internal medicine conference in Wiesbaden, Germany, in 1891. He subsequently published a book on the subject.
The lumbar puncture procedure was taken to the United States by Arthur H. Wentworth an assistant professor at the Harvard Medical School, based at Children's Hospital. In 1893 he published a long paper on diagnosing cerebrospinal meningitis by examining spinal fluid. However, he was criticized by antivivisectionists for having obtained spinal fluid from children. He was acquitted, but, nevertheless, he was uninvited from the then forming Johns Hopkins School of Medicine, where he would have been the first professor of pediatrics.
Historically lumbar punctures were also employed in the process of performing a pneumoencephalography, a nowadays obsolete X-ray imaging study of the brain that was performed extensively from the 1920s until the advent of modern non-invasive neuroimaging techniques such as MRI and CT in the 1970s. During this quite painful procedure, CSF was replaced with air or some other gas via the lumbar puncture in order to enhance the appearance of certain areas of the brain on plain radiographs.
| Biology and health sciences | Medical procedures | null |
342334 | https://en.wikipedia.org/wiki/Autopsy | Autopsy | An autopsy (also referred to as post-mortem examination, obduction, necropsy, or autopsia cadaverum) is a surgical procedure that consists of a thorough examination of a corpse by dissection to determine the cause, mode, and manner of death; or the exam may be performed to evaluate any disease or injury that may be present for research or educational purposes. The term necropsy is generally used for non-human animals.
Autopsies are usually performed by a specialized medical doctor called a pathologist. Only a small portion of deaths require an autopsy to be performed, under certain circumstances. In most cases, a medical examiner or coroner can determine the cause of death.
Purposes of performance
Autopsies are performed for either legal or medical purposes. Autopsies can be performed when any of the following information is desired:
Manner of death must be determined
Determine if death was natural or unnatural
Injury source and extent on the corpse
Post mortem interval
Determining the deceased's identity
Retain relevant organs
If it is an infant, determine live birth and viability
For example, a forensic autopsy is carried out when the cause of death may be a criminal matter, while a clinical or academic autopsy is performed to find the medical cause of death and is used in cases of unknown or uncertain death, or for research purposes. Autopsies can be further classified into cases where an external examination suffices, and those where the body is dissected and an internal examination is conducted. Permission from next of kin may be required for internal autopsy in some cases. Once an internal autopsy is complete, the body is reconstituted by sewing it back together.
Etymology
Autopsy
The term "autopsy" derives from the Ancient Greek αὐτοψία autopsia, "to see for oneself", derived from αὐτός (autos, "oneself") and ὄψις (opsis, "sight, view"). The word has been in use since around the 17th century.
Post-mortem
The term "post-mortem" derives from the Latin post, 'after', and mortem, 'death'. It was first recorded in 1734.
Necropsy
The term "necropsy" is derived from the Greek νεκρός 'death' and ὄψις (opsis, 'sight, view').
Purpose
The principal aims of an autopsy are to determine the cause of death, mode of death, manner of death, the state of health of the person before he or she died, and whether any medical diagnosis and treatment before death were appropriate. In most Western countries the number of autopsies performed in hospitals has been decreasing every year since 1955. Critics, including pathologist and former JAMA editor George D. Lundberg, have charged that the reduction in autopsies is negatively affecting the care delivered in hospitals, because when mistakes result in death, they are often not investigated and lessons, therefore, remain unlearned. When a person has permitted an autopsy in advance of their death, autopsies may also be carried out for the purposes of teaching or medical research. An autopsy is usually performed in cases of sudden death, where a doctor is not able to write a death certificate, or when death is believed to result from an unnatural cause. These examinations are performed under a legal authority (medical examiner, coroner, or procurator fiscal) and do not require the consent of relatives of the deceased. The most extreme example is the examination of murder victims, especially when medical examiners are looking for signs of death or the murder method, such as bullet wounds and exit points, signs of strangulation, or traces of poison.
Some religions including Judaism and Islam usually discourage the performing of autopsies on their adherents. Organizations such as ZAKA in Israel and Misaskim in the United States generally guide families on how to ensure that an unnecessary autopsy is not made.
Autopsies are used in clinical medicine to identify a medical error or a previously unnoticed condition that may endanger the living, such as infectious diseases or exposure to hazardous materials.
A study that focused on myocardial infarction (heart attack) as a cause of death found significant errors of omission and commission, i.e. a sizable number of cases ascribed to myocardial infarctions (MIs) were not MIs and a significant number of non-MIs were MIs.
A systematic review of studies of the autopsy calculated that in about 25% of autopsies, a major diagnostic error will be revealed. However, this rate has decreased over time and the study projects that in a contemporary US institution, 8.4% to 24.4% of autopsies will detect major diagnostic errors.
A large meta-analysis suggested that approximately one-third of death certificates are incorrect and that half of the autopsies performed produced findings that were not suspected before the person died. Also, it is thought that over one-fifth of unexpected findings can only be diagnosed histologically, i.e., by biopsy or autopsy, and that approximately one-quarter of unexpected findings, or 5% of all findings, are major and can similarly only be diagnosed from tissue.
One study found that (out of 694 diagnoses) "Autopsies revealed 171 missed diagnoses, including 21 cancers, 12 strokes, 11 myocardial infarctions, 10 pulmonary emboli, and 9 endocarditis, among others".
Focusing on intubated patients, one study found "abdominal pathologic conditions – abscesses, bowel perforations, or infarction – were as frequent as pulmonary emboli as a cause of class I errors. While patients with abdominal pathologic conditions generally complained of abdominal pain, results of an examination of the abdomen were considered unremarkable in most patients, and the symptom was not pursued".
Types
There are four main types of autopsy:
Medico-legal or forensic or coroner's autopsies seek to find the cause and manner of death and to identify the decedent. They are generally performed, as prescribed by applicable law, in cases of violent, suspicious or sudden deaths, deaths without medical assistance, or during surgical procedures.
Clinical or pathological autopsies are performed to diagnose a particular disease or for research purposes. They aim to determine, clarify, or confirm medical diagnoses that remained unknown or unclear before the patient's death.
Anatomical or academic autopsies are performed by students of anatomy for study purposes only.
Virtual or medical imaging autopsies are performed utilizing imaging technology only, primarily magnetic resonance imaging (MRI) and computed tomography (CT).
Forensic autopsy
A forensic autopsy is used to determine the cause, mode, and manner of death.
Forensic science involves the application of the sciences to answer questions of interest to the legal system.
Medical examiners attempt to determine the time of death, the exact cause of death, and what, if anything, preceded the death, such as a struggle. A forensic autopsy may include obtaining biological specimens from the deceased for toxicological testing, including stomach contents. Toxicology tests may reveal the presence of one or more chemical "poisons" (all chemicals, in sufficient quantities, can be classified as a poison) and their quantity. Because post-mortem deterioration of the body, together with the gravitational pooling of bodily fluids, will necessarily alter the bodily environment, toxicology tests may overestimate, rather than underestimate, the quantity of the suspected chemical.
Following an in-depth examination of all the evidence, a medical examiner or coroner will assign a manner of death from the choices proscribed by the fact-finder's jurisdiction and will detail the evidence on the mechanism of the death.
Clinical autopsy
Clinical autopsies serve two major purposes. They are performed to gain more insight into pathological processes and determine what factors contributed to a patient's death. For example, material for infectious disease testing can be collected during an autopsy. Autopsies are also performed to ensure the standard of care at hospitals. Autopsies can yield insight into how patient deaths can be prevented in the future.
Within the United Kingdom, clinical autopsies can be carried out only with the consent of the family of the deceased person, as opposed to a medico-legal autopsy instructed by a Coroner (England & Wales) or Procurator Fiscal (Scotland), to which the family cannot object.
Over time, autopsies have not only been able to determine the cause of death, but have also led to discoveries of various diseases such as fetal alcohol syndrome, Legionnaire's disease, and even viral hepatitis.
Academic autopsy
Academic autopsies are performed by students of anatomy for the purpose of study, giving medical students and residents firsthand experience viewing anatomy and pathology. Postmortem examinations require the skill to connect anatomic and clinical pathology together since they involve organ systems and interruptions from ante-mortem and post-mortem. These academic autopsies allow for students to practice and develop skills in pathology and become meticulous in later case examinations.
Virtual autopsy
Virtual autopsies are performed using radiographic techniques which can be used in post-mortem examinations for a deceased individual. It is an alternative to medical autopsies, where radiographs are used, for example, Magnetic resonance imaging (MRI) and Computed tomography (CT scan) which produce radiographic images in order to determine the cause of death, the nature, and the manner of death, without dissecting the deceased. It can also be used in the identification of the deceased. This method is helpful in determining the questions pertaining to an autopsy without putting the examiner at risk of biohazardous materials that can be in an individual's body.
Prevalence
In 2004 in England and Wales, there were 514,000 deaths, of which 225,500 were referred to the coroner. Of those, 115,800 (22.5% of all deaths) resulted in post-mortem examinations and there were 28,300 inquests, 570 with a jury.
The rate of consented (hospital) autopsy in the UK and worldwide has declined rapidly over the past 50 years. In the UK in 2013, only 0.7% of inpatient adult deaths were followed by consented autopsy.
The autopsy rate in Germany is below 5% and thus much lower than in other countries in Europe. The governmental reimbursement is hardly sufficient to cover all the costs, so the medical journal Deutsches Ärzteblatt, issued by the German Medical Association, makes the effort to raise awareness regarding the underfinancing of autopsies. The same sources stated that autopsy rates in Sweden and Finland reach 20 to 30%.
In the United States, autopsy rates fell from 17% in 1980 to 14% in 1985 and 11.5% in 1989, although the figures vary notably from county to county.
Process
The body is received at a medical examiner's office, municipal mortuary, or hospital in a body bag or evidence sheet. A new body bag is used for each body to ensure that only evidence from that body is contained within the bag. Evidence sheets are an alternative way to transport the body. An evidence sheet is a sterile sheet that covers the body when it is moved. If it is believed there may be any significant evidence on the hands, for example, gunshot residue or skin under the fingernails, a separate paper sack is put around each hand and taped shut around the wrist.
There are two parts to the physical examination of the body: the external and internal examination. Toxicology, biochemical tests or genetic testing/molecular autopsy often supplement these and frequently assist the pathologist in assigning the cause or causes of death.
External examination
At many institutions, the person responsible for handling, cleaning, and moving the body is called a diener, the German word for servant. In the UK this role is performed by an Anatomical Pathology Technician (APT), who will also assist the pathologist in eviscerating the body and reconstruction after the autopsy. After the body is received, it is first photographed. The examiner then notes the kind of clothes - if any - and their position on the body before they are removed. Next, any evidence such as residue, flakes of paint, or other material is collected from the external surfaces of the body. Ultraviolet light may also be used to search body surfaces for any evidence not easily visible to the naked eye. Samples of hair, nails, and the like are taken, and the body may also be radiographically imaged. Once the external evidence is collected, the body is removed from the bag, undressed, and any wounds present are examined. The body is then cleaned, weighed, and measured in preparation for the internal examination.
A general description of the body as regards ethnic group, sex, age, hair colour and length, eye colour, and other distinguishing features (birthmarks, old scar tissue, moles, tattoos, etc.) is then made. A voice recorder or a standard examination form is normally used to record this information.
In some countries, e.g., Scotland, France, Germany, Russia, and Canada, an autopsy may comprise an external examination only. This concept is sometimes termed a "view and grant". The principle behind this is that the medical records, history of the deceased and circumstances of death have all indicated as to the cause and manner of death without the need for an internal examination.
Internal examination
If not already in place, a plastic or rubber brick called a "head block" is placed under the shoulders of the corpse; hyperflexion of the neck makes the spine arch backward while stretching and pushing the chest upward to make it easier to incise. This gives the APT, or pathologist, maximum exposure to the trunk. After this is done, the internal examination begins. The internal examination consists of inspecting the internal organs of the body by dissection for evidence of trauma or other indications of the cause of death. For the internal examination there are a number of different approaches available:
a large and deep Y-shaped incision can be made starting at the top of each shoulder and running down the front of the chest, meeting at the lower point of the sternum (breastbone).
a curved incision made from the tips of each shoulder, in a semi-circular line across the chest/decolletage, to approximately the level of the second rib, curving back up to the opposite shoulder.
a single vertical incision is made from the sternal notch at the base of the neck.
a U-shaped incision is made at the tip of both shoulders, down along the side of the chest to the bottom of the rib cage, following it. This is typically used on women and during chest-only autopsies.
There is no need for any incision to be made, which will be visible after completion of the examination when the deceased is dressed in a shroud.
In all of the above cases, the incision then extends all the way down to the pubic bone (making a deviation to either side of the navel) and avoiding, where possible, transecting any scars that may be present.
Bleeding from the cuts is minimal, or non-existent because the pull of gravity is producing the only blood pressure at this point, related directly to the complete lack of cardiac functionality. However, in certain cases, there is anecdotal evidence that bleeding can be quite profuse, especially in cases of drowning.
At this point, shears are used to open the chest cavity. The examiner uses the tool to cut through the ribs on the costal cartilage, to allow the sternum to be removed; this is done so that the heart and lungs can be seen in situ and that the heartin particular, the pericardial sacis not damaged or disturbed from opening. A PM 40 knife is used to remove the sternum from the soft tissue that attaches it to the mediastinum. Now the lungs and the heart are exposed. The sternum is set aside and will eventually be replaced at the end of the autopsy.
At this stage, the organs are exposed. Usually, the organs are removed in a systematic fashion. Making a decision as to what order the organs are to be removed will depend highly on the case in question. Organs can be removed in several ways: The first is the en masse technique of Letulle whereby all the organs are removed as one large mass. The second is the en bloc method of Ghon. The most popular in the UK is a modified version of this method, which is divided into four groups of organs. Although these are the two predominant evisceration techniques, in the UK variations on these are widespread.
One method is described here: The pericardial sac is opened to view the heart. Blood for chemical analysis may be removed from the inferior vena cava or the pulmonary veins. Before removing the heart, the pulmonary artery is opened in order to search for a blood clot. The heart can then be removed by cutting the inferior vena cava, the pulmonary veins, the aorta and pulmonary artery, and the superior vena cava. This method leaves the aortic arch intact, which will make things easier for the embalmer. The left lung is then easily accessible and can be removed by cutting the bronchus, artery, and vein at the hilum. The right lung can then be similarly removed. The abdominal organs can be removed one by one after first examining their relationships and vessels.
Most pathologists, however, prefer the organs to be removed all in one "block". Using dissection of the fascia, blunt dissection; using the fingers or hands and traction; the organs are dissected out in one piece for further inspection and sampling. During autopsies of infants, this method is used almost all of the time. The various organs are examined, weighed and tissue samples in the form of slices are taken. Even major blood vessels are cut open and inspected at this stage. Next, the stomach and intestinal contents are examined and weighed. This could be useful to find the cause and time of death, due to the natural passage of food through the bowel during digestion. The more area empty, the longer the deceased had gone without a meal before death.
The body block that was used earlier to elevate the chest cavity is now used to elevate the head. To examine the brain, an incision is made from behind one ear, over the crown of the head, to a point behind the other ear. When the autopsy is completed, the incision can be neatly sewn up and is not noticed when the head is resting on a pillow in an open casket funeral. The scalp is pulled away from the skull in two flaps with the front flap going over the face and the rear flap over the back of the neck. The skull is then cut with a circular (or semicircular) bladed reciprocating saw to create a "cap" that can be pulled off, exposing the brain. The brain is then observed in situ. Then the brain's connections to the cranial nerves and spinal cord are severed, and the brain is lifted out of the skull for further examination. If the brain needs to be preserved before being inspected, it is contained in a large container of formalin (15 percent solution of formaldehyde gas in buffered water) for at least two, but preferably four weeks. This not only preserves the brain, but also makes it firmer, allowing easier handling without corrupting the tissue.
Reconstitution of the body
An important component of the autopsy is the reconstitution of the body such that it can be viewed, if desired, by relatives of the deceased following the procedure. After the examination, the body has an open and empty thoracic cavity with chest flaps open on both sides; the top of the skull is missing, and the skull flaps are pulled over the face and neck. It is unusual to examine the face, arms, hands or legs internally.
In the UK, following the Human Tissue Act 2004 all organs and tissue must be returned to the body unless permission is given by the family to retain any tissue for further investigation. Normally the internal body cavity is lined with cotton, wool, or a similar material, and the organs are then placed into a plastic bag to prevent leakage and are returned to the body cavity. The chest flaps are then closed and sewn back together and the skull cap is sewed back in place. Then the body may be wrapped in a shroud, and it is common for relatives to not be able to tell the procedure has been done when the body is viewed in a funeral parlor after embalming.
In stroke
An autopsy of stroke may be able to establish the time taken from the onset of cerebral infarction to the time of death.
Various microscopic findings are present at times from infarction as follows:
History
Around 3000 BCE, ancient Egyptians were one of the first civilizations to practice the removal and examination of the internal organs of humans in the religious practice of mummification.
Autopsies that opened the body to determine the cause of death were attested at least in the early third millennium BCE, although they were opposed in many ancient societies where it was believed that the outward disfigurement of dead persons prevented them from entering the afterlife (as with the Egyptians, who removed the organs through tiny slits in the body). Notable Greek autopsists were Erasistratus and Herophilus of Chalcedon, who lived in 3rd century BCE Alexandria, but in general, autopsies were rare in ancient Greece. In 44 BCE, Julius Caesar was the subject of an official autopsy after his murder by rival senators, the physician's report noting that the second stab wound Caesar received was the fatal one. Julius Caesar had been stabbed a total of 23 times. By around 150 BCE, ancient Roman legal practice had established clear parameters for autopsies. The greatest ancient anatomist was Galen (CE 129– ), whose findings would not be challenged until the Renaissance over a thousand years later.
Ibn Tufail has elaborated on autopsy in his treatise called Hayy ibn Yaqzan and Nadia Maftouni, discussing the subject in an extensive article, believes him to be among the early supporters of autopsy and vivisection.
The dissection of human remains for medical or scientific reasons continued to be practiced irregularly after the Romans, for instance by the Arab physicians Avenzoar and Ibn al-Nafis. In Europe they were done with enough regularity to become skilled, as early as 1200, and successful efforts to preserve the body, by filling the veins with wax and metals. Until the 20th century, it was thought that the modern autopsy process derived from the anatomists of the Renaissance. Giovanni Battista Morgagni (1682–1771), celebrated as the father of anatomical pathology, wrote the first exhaustive work on pathology, De Sedibus et Causis Morborum per Anatomen Indagatis (The Seats and Causes of Diseases Investigated by Anatomy, 1769).
In 1543, Andreas Vesalius conducted a public dissection of the body of a former criminal. He asserted and articulated the bones, this became the world's oldest surviving anatomical preparation. It is still displayed at the Anatomical Museum at the University of Basel.
In the mid-1800s, Carl von Rokitansky and colleagues at the Second Vienna Medical School began to undertake dissections as a means to improve diagnostic medicine.
The 19th-century medical researcher Rudolf Virchow, in response to a lack of standardization of autopsy procedures, established and published specific autopsy protocols (one such protocol still bears his name). He also developed the concept of pathological processes.
During the turn of the 20th century, the Scotland Yard created the Office of the Forensic Pathologist, a medical examiner trained in medicine, charged with investigating the cause of all unnatural deaths, including accidents, homicides, suicides, etc.
Other animals (necropsy)
A post-mortem examination, or necropsy, is far more common in veterinary medicine than in human medicine. For many species that exhibit few external symptoms (sheep), or that are not suited to detailed clinical examination (poultry, cage birds, zoo animals), it is a common method used by veterinary physicians to come to a diagnosis. A necropsy is mostly used like an autopsy to determine the cause of death. The entire body is examined at the gross visual level, and samples are collected for additional analyses.
| Biology and health sciences | Medical procedures | null |
342684 | https://en.wikipedia.org/wiki/Inclusion%E2%80%93exclusion%20principle | Inclusion–exclusion principle | In combinatorics, the inclusion–exclusion principle is a counting technique which generalizes the familiar method of obtaining the number of elements in the union of two finite sets; symbolically expressed as
where A and B are two finite sets and |S| indicates the cardinality of a set S (which may be considered as the number of elements of the set, if the set is finite). The formula expresses the fact that the sum of the sizes of the two sets may be too large since some elements may be counted twice. The double-counted elements are those in the intersection of the two sets and the count is corrected by subtracting the size of the intersection.
The inclusion-exclusion principle, being a generalization of the two-set case, is perhaps more clearly seen in the case of three sets, which for the sets A, B and C is given by
This formula can be verified by counting how many times each region in the Venn diagram figure is included in the right-hand side of the formula. In this case, when removing the contributions of over-counted elements, the number of elements in the mutual intersection of the three sets has been subtracted too often, so must be added back in to get the correct total.
Generalizing the results of these examples gives the principle of inclusion–exclusion. To find the cardinality of the union of sets:
Include the cardinalities of the sets.
Exclude the cardinalities of the pairwise intersections.
Include the cardinalities of the triple-wise intersections.
Exclude the cardinalities of the quadruple-wise intersections.
Include the cardinalities of the quintuple-wise intersections.
Continue, until the cardinality of the -tuple-wise intersection is included (if is odd) or excluded ( even).
The name comes from the idea that the principle is based on over-generous inclusion, followed by compensating exclusion.
This concept is attributed to Abraham de Moivre (1718), although it first appears in a paper of Daniel da Silva (1854) and later in a paper by J. J. Sylvester (1883). Sometimes the principle is referred to as the formula of Da Silva or Sylvester, due to these publications. The principle can be viewed as an example of the sieve method extensively used in number theory and is sometimes referred to as the sieve formula.
As finite probabilities are computed as counts relative to the cardinality of the probability space, the formulas for the principle of inclusion–exclusion remain valid when the cardinalities of the sets are replaced by finite probabilities. More generally, both versions of the principle can be put under the common umbrella of measure theory.
In a very abstract setting, the principle of inclusion–exclusion can be expressed as the calculation of the inverse of a certain matrix. This inverse has a special structure, making the principle an extremely valuable technique in combinatorics and related areas of mathematics. As Gian-Carlo Rota put it:
"One of the most useful principles of enumeration in discrete probability and combinatorial theory is the celebrated principle of inclusion–exclusion. When skillfully applied, this principle has yielded the solution to many a combinatorial problem."
Formula
In its general formula, the principle of inclusion–exclusion states that for finite sets , one has the identity
This can be compactly written as
or
In words, to count the number of elements in a finite union of finite sets, first sum the cardinalities of the individual sets, then subtract the number of elements that appear in at least two sets, then add back the number of elements that appear in at least three sets, then subtract the number of elements that appear in at least four sets, and so on. This process always ends since there can be no elements that appear in more than the number of sets in the union. (For example, if there can be no elements that appear in more than sets; equivalently, there can be no elements that appear in at least sets.)
In applications it is common to see the principle expressed in its complementary form. That is, letting be a finite universal set containing all of the and letting denote the complement of in , by De Morgan's laws we have
As another variant of the statement, let be a list of properties that elements of a set may or may not have, then the principle of inclusion–exclusion provides a way to calculate the number of elements of that have none of the properties. Just let be the subset of elements of which have the property and use the principle in its complementary form. This variant is due to J. J. Sylvester.
Notice that if you take into account only the first sums on the right (in the general form of the principle), then you will get an overestimate if is odd and an underestimate if is even.
Examples
Counting derangements
A more complex example is the following.
Suppose there is a deck of n cards numbered from 1 to n. Suppose a card numbered m is in the correct position if it is the mth card in the deck. How many ways, W, can the cards be shuffled with at least 1 card being in the correct position?
Begin by defining set Am, which is all of the orderings of cards with the mth card correct. Then the number of orders, W, with at least one card being in the correct position, m, is
Apply the principle of inclusion–exclusion,
Each value represents the set of shuffles having at least p values m1, ..., mp in the correct position. Note that the number of shuffles with at least p values correct only depends on p, not on the particular values of . For example, the number of shuffles having the 1st, 3rd, and 17th cards in the correct position is the same as the number of shuffles having the 2nd, 5th, and 13th cards in the correct positions. It only matters that of the n cards, 3 were chosen to be in the correct position. Thus there are equal terms in the pth summation (see combination).
is the number of orderings having p elements in the correct position, which is equal to the number of ways of ordering the remaining n − p elements, or (n − p)!. Thus we finally get:
A permutation where no card is in the correct position is called a derangement. Taking n! to be the total number of permutations, the probability Q that a random shuffle produces a derangement is given by
a truncation to n + 1 terms of the Taylor expansion of e−1. Thus the probability of guessing an order for a shuffled deck of cards and being incorrect about every card is approximately e−1 or 37%.
A special case
The situation that appears in the derangement example above occurs often enough to merit special attention. Namely, when the size of the intersection sets appearing in the formulas for the principle of inclusion–exclusion depend only on the number of sets in the intersections and not on which sets appear. More formally, if the intersection
has the same cardinality, say αk = |AJ|, for every k-element subset J of {1, ..., n}, then
Or, in the complementary form, where the universal set S has cardinality α0,
Formula generalization
Given a family (repeats allowed) of subsets A1, A2, ..., An of a universal set S, the principle of inclusion–exclusion calculates the number of elements of S in none of these subsets. A generalization of this concept would calculate the number of elements of S which appear in exactly some fixed m of these sets.
Let N = [n] = {1,2,...,n}. If we define , then the principle of inclusion–exclusion can be written as, using the notation of the previous section; the number of elements of S contained in none of the Ai is:
If I is a fixed subset of the index set N, then the number of elements which belong to Ai for all i in I and for no other values is:
Define the sets
We seek the number of elements in none of the Bk which, by the principle of inclusion–exclusion (with ), is
The correspondence K ↔ J = I ∪ K between subsets of N \ I and subsets of N containing I is a bijection and if J and K correspond under this map then BK = AJ, showing that the result is valid.
In probability
In probability, for events A1, ..., An in a probability space , the inclusion–exclusion principle becomes for n = 2
for n = 3
and in general
which can be written in closed form as
where the last sum runs over all subsets I of the indices 1, ..., n which contain exactly k elements, and
denotes the intersection of all those Ai with index in I.
According to the Bonferroni inequalities, the sum of the first terms in the formula is alternately an upper bound and a lower bound for the LHS. This can be used in cases where the full formula is too cumbersome.
For a general measure space (S,Σ,μ) and measurable subsets A1, ..., An of finite measure, the above identities also hold when the probability measure is replaced by the measure μ.
Special case
If, in the probabilistic version of the inclusion–exclusion principle, the probability of the intersection AI only depends on the cardinality of I, meaning that for every k in {1, ..., n} there is an ak such that
then the above formula simplifies to
due to the combinatorial interpretation of the binomial coefficient . For example, if the events are independent and identically distributed, then for all i, and we have , in which case the expression above simplifies to
(This result can also be derived more simply by considering the intersection of the complements of the events .)
An analogous simplification is possible in the case of a general measure space and measurable subsets of finite measure.
There is another formula used in point processes. Let be a finite set and be a random subset of . Let be any subset of , then
Other formulas
The principle is sometimes stated in the form that says that if
then
The combinatorial and the probabilistic version of the inclusion–exclusion principle are instances of ().
If one sees a number as a set of its prime factors, then () is a generalization of Möbius inversion formula for square-free natural numbers. Therefore, () is seen as the Möbius inversion formula for the incidence algebra of the partially ordered set of all subsets of A.
For a generalization of the full version of Möbius inversion formula, () must be generalized to multisets. For multisets instead of sets, () becomes
where is the multiset for which , and
μ(S) = 1 if S is a set (i.e. a multiset without double elements) of even cardinality.
μ(S) = −1 if S is a set (i.e. a multiset without double elements) of odd cardinality.
μ(S) = 0 if S is a proper multiset (i.e. S has double elements).
Notice that is just the of () in case is a set.
Applications
The inclusion–exclusion principle is widely used and only a few of its applications can be mentioned here.
Counting derangements
A well-known application of the inclusion–exclusion principle is to the combinatorial problem of counting all derangements of a finite set. A derangement of a set A is a bijection from A into itself that has no fixed points. Via the inclusion–exclusion principle one can show that if the cardinality of A is n, then the number of derangements is [n! / e] where [x] denotes the nearest integer to x; a detailed proof is available here and also see the examples section above.
The first occurrence of the problem of counting the number of derangements is in an early book on games of chance: Essai d'analyse sur les jeux de hazard by P. R. de Montmort (1678 – 1719) and was known as either "Montmort's problem" or by the name he gave it, "problème des rencontres." The problem is also known as the hatcheck problem.
The number of derangements is also known as the subfactorial of n, written !n. It follows that if all bijections are assigned the same probability then the probability that a random bijection is a derangement quickly approaches 1/e as n grows.
Counting intersections
The principle of inclusion–exclusion, combined with De Morgan's law, can be used to count the cardinality of the intersection of sets as well. Let represent the complement of Ak with respect to some universal set A such that for each k. Then we have
thereby turning the problem of finding an intersection into the problem of finding a union.
Graph coloring
The inclusion exclusion principle forms the basis of algorithms for a number of NP-hard graph partitioning problems, such as graph coloring.
A well known application of the principle is the construction of the chromatic polynomial of a graph.
Bipartite graph perfect matchings
The number of perfect matchings of a bipartite graph can be calculated using the principle.
Number of onto functions
Given finite sets A and B, how many surjective functions (onto functions) are there from A to B? Without any loss of generality we may take A = {1, ..., k} and B = {1, ..., n}, since only the cardinalities of the sets matter. By using S as the set of all functions from A to B, and defining, for each i in B, the property Pi as "the function misses the element i in B" (i is not in the image of the function), the principle of inclusion–exclusion gives the number of onto functions between A and B as:
Permutations with forbidden positions
A permutation of the set S = {1, ..., n} where each element of S is restricted to not being in certain positions (here the permutation is considered as an ordering of the elements of S) is called a permutation with forbidden positions. For example, with S = {1,2,3,4}, the permutations with the restriction that the element 1 can not be in positions 1 or 3, and the element 2 can not be in position 4 are: 2134, 2143, 3124, 4123, 2341, 2431, 3241, 3421, 4231 and 4321. By letting Ai be the set of positions that the element i is not allowed to be in, and the property Pi to be the property that a permutation puts element i into a position in Ai, the principle of inclusion–exclusion can be used to count the number of permutations which satisfy all the restrictions.
In the given example, there are 12 = 2(3!) permutations with property P1, 6 = 3! permutations with property P2 and no permutations have properties P3 or P4 as there are no restrictions for these two elements. The number of permutations satisfying the restrictions is thus:
4! − (12 + 6 + 0 + 0) + (4) = 24 − 18 + 4 = 10.
The final 4 in this computation is the number of permutations having both properties P1 and P2. There are no other non-zero contributions to the formula.
Stirling numbers of the second kind
The Stirling numbers of the second kind, S(n,k) count the number of partitions of a set of n elements into k non-empty subsets (indistinguishable boxes). An explicit formula for them can be obtained by applying the principle of inclusion–exclusion to a very closely related problem, namely, counting the number of partitions of an n-set into k non-empty but distinguishable boxes (ordered non-empty subsets). Using the universal set consisting of all partitions of the n-set into k (possibly empty) distinguishable boxes, A1, A2, ..., Ak, and the properties Pi meaning that the partition has box Ai empty, the principle of inclusion–exclusion gives an answer for the related result. Dividing by k! to remove the artificial ordering gives the Stirling number of the second kind:
Rook polynomials
A rook polynomial is the generating function of the number of ways to place non-attacking rooks on a board B that looks like a subset of the squares of a checkerboard; that is, no two rooks may be in the same row or column. The board B is any subset of the squares of a rectangular board with n rows and m columns; we think of it as the squares in which one is allowed to put a rook. The coefficient, rk(B) of xk in the rook polynomial RB(x) is the number of ways k rooks, none of which attacks another, can be arranged in the squares of B. For any board B, there is a complementary board consisting of the squares of the rectangular board that are not in B. This complementary board also has a rook polynomial with coefficients
It is sometimes convenient to be able to calculate the highest coefficient of a rook polynomial in terms of the coefficients of the rook polynomial of the complementary board. Without loss of generality we can assume that n ≤ m, so this coefficient is rn(B). The number of ways to place n non-attacking rooks on the complete n × m "checkerboard" (without regard as to whether the rooks are placed in the squares of the board B) is given by the falling factorial:
Letting Pi be the property that an assignment of n non-attacking rooks on the complete board has a rook in column i which is not in a square of the board B, then by the principle of inclusion–exclusion we have:
Euler's phi function
Euler's totient or phi function, φ(n) is an arithmetic function that counts the number of positive integers less than or equal to n that are relatively prime to n. That is, if n is a positive integer, then φ(n) is the number of integers k in the range 1 ≤ k ≤ n which have no common factor with n other than 1. The principle of inclusion–exclusion is used to obtain a formula for φ(n). Let S be the set {1, ..., n} and define the property Pi to be that a number in S is divisible by the prime number pi, for 1 ≤ i ≤ r, where the prime factorization of
Then,
Dirichlet hyperbola method
The Dirichlet hyperbola method re-expresses a sum of a multiplicative function by selecting a suitable Dirichlet convolution , recognizing that the sum
can be recast as a sum over the lattice points in a region bounded by , , and , splitting this region into two overlapping subregions, and finally using the inclusion–exclusion principle to conclude that
Diluted inclusion–exclusion principle
In many cases where the principle could give an exact formula (in particular, counting prime numbers using the sieve of Eratosthenes), the formula arising does not offer useful content because the number of terms in it is excessive. If each term individually can be estimated accurately, the accumulation of errors may imply that the inclusion–exclusion formula is not directly applicable. In number theory, this difficulty was addressed by Viggo Brun. After a slow start, his ideas were taken up by others, and a large variety of sieve methods developed. These for example may try to find upper bounds for the "sieved" sets, rather than an exact formula.
Let A1, ..., An be arbitrary sets and p1, ..., pn real numbers in the closed unit interval . Then, for every even number k in {0, ..., n}, the indicator functions satisfy the inequality:
Proof of main statement
Choose an element contained in the union of all sets and let be the individual sets containing it. (Note that t > 0.) Since the element is counted precisely once by the left-hand side of equation (), we need to show that it is counted precisely once by the right-hand side. On the right-hand side, the only non-zero contributions occur when all the subsets in a particular term contain the chosen element, that is, all the subsets are selected from . The contribution is one for each of these sets (plus or minus depending on the term) and therefore is just the (signed) number of these subsets used in the term. We then have:
By the binomial theorem,
Using the fact that and rearranging terms, we have
and so, the chosen element is counted only once by the right-hand side of equation ().
Algebraic proof
An algebraic proof can be obtained using indicator functions (also known as characteristic functions). The indicator function of a subset S of a set X is the function
If and are two subsets of , then
Let A denote the union of the sets A1, ..., An. To prove the inclusion–exclusion principle in general, we first verify the identity
for indicator functions, where:
The following function
is identically zero because: if x is not in A, then all factors are 0−0 = 0; and otherwise, if x does belong to some Am, then the corresponding mth factor is 1−1=0. By expanding the product on the left-hand side, equation () follows.
To prove the inclusion–exclusion principle for the cardinality of sets, sum the equation () over all x in the union of A1, ..., An. To derive the version used in probability, take the expectation in (). In general, integrate the equation () with respect to μ. Always use linearity in these derivations.
| Mathematics | Combinatorics | null |
342941 | https://en.wikipedia.org/wiki/Maritime%20patrol%20aircraft | Maritime patrol aircraft | A maritime patrol aircraft (MPA), also known as a patrol aircraft, maritime reconnaissance aircraft, maritime surveillance aircraft, or by the older American term patrol bomber, is a fixed-wing aircraft designed to operate for long durations over water in maritime patrol roles — in particular anti-submarine warfare (ASW), anti-ship warfare (AShW), and search and rescue (SAR).
Among other maritime surveillance resources, such as satellites, ships, unmanned aerial vehicles (UAVs) and helicopters, the MPA is an important asset. To perform ASW operations, MPAs typically carry air-deployable sonar buoys as well as torpedoes and are usually capable of extended flight at low altitudes.
History
First World War
The first aircraft that would now be identified as maritime patrol aircraft were flown by the Royal Naval Air Service and the French Aéronautique Maritime during the First World War, primarily on anti-submarine patrols. France, Italy and Austria-Hungary used large numbers of smaller patrol aircraft for the Mediterranean, Adriatic and other coastal areas while the Germans and British fought over the North Sea. At first, blimps and zeppelins were the only aircraft capable of staying aloft for the longer ten hour patrols whilst carrying a useful payload while shorter-range patrols were mounted with landplanes such as the Sopwith 1½ Strutter. A number of specialized patrol balloons were built, particularly by the British, including the SS class airship of which 158 were built including subtypes.
As the conflict continued, numerous aircraft were developed specifically for the role, including small flying boats such as the FBA Type C, as well as large floatplanes such as the Short 184, or flying boats such as the Felixstowe F.3. Developments of the Felixstowe served with the Royal Air Force until the mid 20s, and with the US Navy as the Curtiss F5L and Naval Aircraft Factory PN whose developments saw service until 1938. During the war, Dornier did considerable pioneering work in all aluminium aircraft structures while working for Luftschiffbau Zeppelin and built four large patrol flying boats, the last of which, the Zeppelin-Lindau Rs.IV, influenced development elsewhere resulting in the replacement of wooden hulls with metal ones, such as on the Short Singapore. The success of long range patrol aircraft led to the development of fighters specifically designed to intercept them, such as the Hansa-Brandenburg W.29.
Second World War
Many of the Second World War patrol airplanes were converted from either bombers or airliners, such as the Lockheed Hudson which started out as the Lockheed Model 14 Super Electra, as well as older biplane designs such as the Supermarine Stranraer, which had begun to be replaced by monoplanes just before the outbreak of war. The British in particular used obsolete bombers to supplement purpose-built aircraft for maritime patrol, such as the Vickers Wellington and Armstrong-Whitworth Whitley, while the US relegated the Douglas B-18 Bolo to the same role until better aircraft became available. Blimps were widely used by the United States Navy, especially in the warmer and calmer latitudes of the Caribbean Sea, the Bahamas, Bermuda, the Gulf of Mexico, Puerto Rico, Trinidad, and later the Azores.
A number of special-purpose aircraft were also used in the conflict, including the American-made twin-engine Consolidated PBY Catalina flying boats, and the large, four-engine British Short Sunderland flying boats of the Allies. In the Pacific theatre, the Catalina was gradually superseded by the longer-ranged Martin PBM Mariner flying boat. For the Axis Powers, there were the long-range Japanese Kawanishi H6K and Kawanishi H8K flying boats, and the German Blohm & Voss BV 138 diesel-engined trimotor flying boat, as well as the converted Focke-Wulf Fw 200 Condor airliner landplane.
To finally close the Mid-Atlantic gap, or "Black Gap", a space in which Axis submarines could prey on Allied shipping out of reach of MPAs, the British Royal Air Force, the Royal Canadian Air Force, and the US Army Air Forces introduced the American Consolidated B-24 Liberator bomber, which had a very long range for the era. The B-24 was also used at the basis for the PB4Y-2 Privateer, a dedicated MPA variant adopted in large numbers by the US Navy, which saw service late on in the Pacific theatre.
During the conflict, there were several developments in air-to-surface-vessel radar and sonobuoys, which enhanced the ability of aircraft to find and destroy submarines, especially at night and in poor weather. Another area of advancement was the adoption of increasingly effective camouflage schemes, which led to the widespread adoption of white paint schemes in the Atlantic to reduce the warning available to surfaced U-boats, while US Navy aircraft transitioned from an upper light blue-gray and lower white to an all-over dark blue due to the increasing threat of Japanese forces at night-time.
Cold War era
In the decades following the Second World War, the MPA missions were partially taken over by aircraft derived from civilian airliners. These had range and performance factors better than most of the wartime bombers. The latest jet-powered bombers of the 1950s did not have the endurance needed for long, overwater patrolling, and they did not have the low loitering speeds necessary for antisubmarine operations. The main threat to NATO maritime supremacy throughout the 1960s, 1970s, and the 1980s was Soviet Navy and Warsaw Pact submarines. These were countered by the NATO fleets, the NATO patrol planes mentioned above, and by sophisticated underwater listening systems. These span the so-called "GIUK Gap" of the North Atlantic that extends from Greenland to Iceland, to the Faroe Islands, to Scotland in the United Kingdom. Air bases for NATO patrol planes have also been located in these areas: U.S. Navy and Canadian aircraft based in Greenland, Iceland, and Newfoundland; British aircraft based in Scotland and Northern Ireland; and Norwegian, Dutch, and German aircraft based in their home countries.
During the late 1940s, the RAF introduced the Avro Shackleton a specialised MPA derivative of the Avro Lancaster bomber in anticipation of a rapid expansion of the Soviet Navy's submarine force.
An improved model of the Shackleton, the MR 3, was introduced, featuring various structural improvements, along with homing torpedoes and Mk 101 Lulu nuclear depth bombs. During the late 1960s, a jet-powered replacement in the form of the Hawker Siddeley Nimrod, a derivation of the De Havilland Comet airliner, begun to be introduced. During the 2000s, an improved model, the BAE Systems Nimrod MRA4, was in development, but was cancelled and eventually substituted for by the Boeing P-8 Poseidon.
The U.S. Navy flew a mixture of MPAs, including the land-based Lockheed P2V Neptune (P2V) and the carrier-based Grumman S-2 Tracker. During the 1970s, the P2V was entirely replaced by the Lockheed P-3 Orion, which remained in service into the early twenty-first century.
The P-3, powered by four turboprop engines, is derived from the 1950s era Lockheed Electra airliner. In addition to their ASW and SAR capabilities, most P-3Cs have been modified to carry Harpoon and Maverick missiles for attacking surface ships. American P-3s were formerly armed with the Lulu nuclear depth charge for ASW, but those were removed from the arsenal and scrapped decades ago. Produced in United States, Japan and Canada, the P-3 has been operated by the air forces and navies of United States, Japan, Canada, Australia, Iran, Brazil, Germany, the Netherlands, New Zealand, Norway, Spain, and Taiwan. The Canadian version is called the CP-140 Aurora.
During the 1960s, in response to North Atlantic Treaty Organization (NATO) issuing a Request for Proposals (RFP) for a new MPA, the Breguet 1150 Atlantic was developed by a French-led multinational consortium, Société d'Étude et de Construction de Breguet Atlantic (SECBAT). Operators of the type include the French Navy, the German Navy, the Italian Air Force, the Pakistan Navy, and the Royal Netherlands Navy. During the 1980s, an updated version, the Atlantic Nouvelle Génération or Atlantique 2, with new equipment and avionics was introduced, which included a new radar, sonar processor, forward-looking infrared camera turret, and the ability to carry the Exocet anti-shipping missile. By 2005, French manufacturer Dassault Aviation had decided to terminate marketing efforts for the Atlantic, promoting a MPA variant of the Dassault Falcon 900 corporate jet instead.
Japan has developed multiple purpose-designed MPAs during this period. The Shin Meiwa PS-1 flying boat was designed to meet a Japanese requirement for a new ASW platform. A modernised derivative of the PS-1, the ShinMaywa US-2 amphibian, was introduced during the early twenty-first century to succeed the PS-1. The land-based Kawasaki P-1 was introduced during the 2010s by the Japan Maritime Self-Defense Force (JMSDF) as a replacement for the aging P-3C Orion.
Both the Royal Australian Air Force and the Royal Australian Navy met their early postwar MPA needs via a stretched-fuselage modification of the Avro Lincoln bomber. However, the type was soon supplemented and eventually replaced by new aircraft, such as the P2V and later the P-3C, which later became the sole ASW type operated by the service.
The Soviet Union developed the Ilyushin Il-38 from a civilian airliner. Similarly, the Royal Canadian Air Force derived the Canadair CP-107 Argus from a British airliner, the Bristol Britannia. The Argus was superseded by the CP-140 Aurora, derived from the Lockheed Electra.
Since the end of the Cold War, the threat of a large-scale submarine attack is a remote one, and many of the air forces and navies have been downsizing their fleets of patrol planes. Those still in service are still used for search-and-rescue, counter-smuggling, antipiracy, antipoaching of marine life, the enforcement of the exclusive economic zones, and enforcement of the laws of the seas.
Armament and countermeasures
The earliest patrol aircraft carried bombs and machine guns. Between the wars the British experimented with equipping their patrol aircraft with the COW 37 mm gun. During World War II, depth charges that could be set to detonate at specific depths, and later when in proximity with large metal objects replaced anti-submarine bombs that detonated on contact. Patrol aircraft also carried defensive armament which was necessary when patrolling areas close to enemy territory such as Allied operations in the Bay of Biscay targeting U-boats starting out from their base.
As a result of Allied successes with patrol aircraft against U-boats, the Germans introduced U-flak (submarines equipped with more antiaircraft weaponry) to escort U-boats out of base and encouraged commanders to remain on the surface and fire back at attacking craft rather than trying to escape by diving. However, U-flak was short-lived, as opposing pilots adapted their tactics. Equipping submarines with radar warning receivers and the snorkel made them harder to find.
To counter the German long-range patrol aircraft that targeted merchant convoys, the Royal Navy introduced the "CAM ship", which was a merchant vessel equipped with a lone fighter plane which could be launched once to engage the enemy planes. Later, the small escort carriers of WWII became available to cover the deep oceans, and the land air bases in the Azores became available in mid-1943 from Portugal.
As technology progressed the bombs and depth charges were supplemented with Acoustic torpedoes that could detect, follow and then explode against an enemy submarine. The US Navy began fielding the Mark 24 mine in 1943, labelled as a mine as a security measure. It sank 37 Axis submarines during the war.
The Cold War era saw the introduction of the nuclear depth bomb, a depth charge with a nuclear warhead that raised the probability of a kill against a submarine to a near-certainty as long as detonation occurred.
While anti-submarine warfare is the main role of patrol aircraft, their large payload capability has seen them fitted for various weaponry outside their nominal role. The Lockheed P-3 Orion was fitted with underwing pylons that could carry a variety of common American weapons, including the AGM-84 Harpoon anti-ship missile, the air-to-ground AGM-65 Maverick, as many as ten of the CBU-100 Cluster Bomb, rocket pods, sea mines, and the standard issue Mark 80 general purpose bombs.
The Royal Air Force's Hawker Siddeley Nimrod was fitted with AIM-9 Sidewinder missiles in 1982 during the Falklands War in order for it to be able to attack any Argentine Air Force patrol planes they might encounter.
Sensors
Maritime patrol aircraft are typically fitted with a wide range of sensors:
Radar to detect surface shipping movements. Radar can also detect a submarine snorkel or periscope, and the wake it creates.
Magnetic anomaly detector (MAD) to detect the iron in a submarine's hull. The MAD sensor is typically mounted on an extension from the tail or is trailed behind the aircraft on a cable to minimize interference from the metal in the rest of the aircraft;
Sonobuoys - self-contained sonar transmitter/receivers dropped into the water to transmit data back to the aircraft for analysis;
ELINT sensors to monitor communications and radar emissions;
Infrared cameras (sometimes referred to as FLIR for forward looking infrared) for detecting exhaust streams and other sources of heat and are useful in monitoring shipping movements and fishing activity.
Visual inspection using the aircrew's eyes, in some cases aided by searchlights or flares.
A modern military maritime patrol aircraft typically carries a dozen or so crew members, including relief flight crews, to effectively operate the equipment for 12 hours or more at a time.
Examples
| Technology | Military aviation | null |
343225 | https://en.wikipedia.org/wiki/Sea%20ice | Sea ice | Sea ice arises as seawater freezes. Because ice is less dense than water, it floats on the ocean's surface (as does fresh water ice). Sea ice covers about 7% of the Earth's surface and about 12% of the world's oceans. Much of the world's sea ice is enclosed within the polar ice packs in the Earth's polar regions: the Arctic ice pack of the Arctic Ocean and the Antarctic ice pack of the Southern Ocean. Polar packs undergo a significant yearly cycling in surface extent, a natural process upon which depends the Arctic ecology, including the ocean's ecosystems. Due to the action of winds, currents and temperature fluctuations, sea ice is very dynamic, leading to a wide variety of ice types and features. Sea ice may be contrasted with icebergs, which are chunks of ice shelves or glaciers that calve into the ocean. Depending on location, sea ice expanses may also incorporate icebergs.
General features and dynamics
Sea ice does not simply grow and melt. During its lifespan, it is very dynamic. Due to the combined action of winds, currents, water temperature and air temperature fluctuations, sea ice expanses typically undergo a significant amount of deformation. Sea ice is classified according to whether or not it is able to drift and according to its age.
Fast ice versus drift (or pack) ice
Sea ice can be classified according to whether or not it is attached (or frozen) to the shoreline (or between shoals or to grounded icebergs). If attached, it is called landfast ice, or more often, fast ice (as in fastened). Alternatively and unlike fast ice, drift ice occurs further offshore in very wide areas and encompasses ice that is free to move with currents and winds. The physical boundary between fast ice and drift ice is the fast ice boundary. The drift ice zone may be further divided into a shear zone, a marginal ice zone and a central pack. Drift ice consists of floes, individual pieces of sea ice or more across. There are names for various floe sizes: small – ; medium – ; big – ; vast – ; and giant – more than . The term pack ice is used either as a synonym to drift ice, or to designate drift ice zone in which the floes are densely packed. The overall sea ice cover is termed the ice canopy from the perspective of submarine navigation.
Classification based on age
Another classification used by scientists to describe sea ice is based on age, that is, on its development stages. These stages are: new ice, nilas, young ice, first-year and old.
New ice, nilas and young ice
New ice is a general term used for recently frozen sea water that does not yet make up solid ice. It may consist of frazil ice (plates or spicules of ice suspended in water), slush (water saturated snow), or shuga (spongy white ice lumps a few centimeters across). Other terms, such as grease ice and pancake ice, are used for ice crystal accumulations under the action of wind and waves. When sea ice begins to form on a beach with a light swell, ice eggs up to the size of a football can be created.
Nilas designates a sea ice crust up to in thickness. It bends without breaking around waves and swells. Nilas can be further subdivided into dark nilas – up to in thickness and very dark and light nilas – over in thickness and lighter in color.
Young ice is a transition stage between nilas and first-year ice and ranges in thickness from to , Young ice can be further subdivided into grey ice – to in thickness and grey-white ice – to in thickness. Young ice is not as flexible as nilas, but tends to break under wave action. Under compression, it will either raft (at the grey ice stage) or ridge (at the grey-white ice stage).
First-year sea ice
First-year sea ice is ice that is thicker than young ice but has no more than one year growth. In other words, it is ice that grows in the fall and winter (after it has gone through the new ice – nilas – young ice stages and grows further) but does not survive the spring and summer months (it melts away). The thickness of this ice typically ranges from to . First-year ice may be further divided into thin ( to ), medium ( to ) and thick (>).
Old sea ice
Old sea ice is sea ice that has survived at least one melting season (i.e. one summer). For this reason, this ice is generally thicker than first-year sea ice. The thickness of old sea ice typically ranges from 2 to 4 m. Old ice is commonly divided into two types: second-year ice, which has survived one melting season and multiyear ice, which has survived more than one. (In some sources, old ice is more than two years old.) Multi-year ice is much more common in the Arctic than it is in the Antarctic. The reason for this is that sea ice in the south drifts into warmer waters where it melts. In the Arctic, much of the sea ice is land-locked.
Driving forces
While fast ice is relatively stable (because it is attached to the shoreline or the seabed), drift (or pack) ice undergoes relatively complex deformation processes that ultimately give rise to sea ice's typically wide variety of landscapes. Wind is the main driving force, along with ocean currents. The Coriolis force and sea ice surface tilt have also been invoked. These driving forces induce a state of stress within the drift ice zone. An ice floe converging toward another and pushing against it will generate a state of compression at the boundary between both. The ice cover may also undergo a state of tension, resulting in divergence and fissure opening. If two floes drift sideways past each other while remaining in contact, this will create a state of shear.
Deformation
Sea ice deformation results from the interaction between ice floes as they are driven against each other. The result may be of three types of features: 1) Rafted ice, when one piece is overriding another; 2) Pressure ridges, a line of broken ice forced downward (to make up the keel) and upward (to make the sail); and 3) Hummock, a hillock of broken ice that forms an uneven surface. A shear ridge is a pressure ridge that formed under shear – it tends to be more linear than a ridge induced only by compression. A new ridge is a recent feature – it is sharp-crested, with its side sloping at an angle exceeding 40 degrees. In contrast, a weathered ridge is one with a rounded crest and with sides sloping at less than 40 degrees. Stamukhi are yet another type of pile-up but these are grounded and are therefore relatively stationary. They result from the interaction between fast ice and the drifting pack ice.
Level ice is sea ice that has not been affected by deformation and is therefore relatively flat.
Leads and polynyas
Leads and polynyas are areas of open water that occur within sea ice expanses even though air temperatures are below freezing. They provide a direct interaction between the ocean and the atmosphere, which is important for the wildlife. Leads are narrow and linear, varying in width from meters to kilometers. During the winter, the water in leads quickly freezes up. They are also used for navigation purposes – even when refrozen, the ice in leads is thinner, allowing icebreakers access to an easier sail path and submarines to surface more easily. Polynyas are more uniform in size than leads and are also larger – two types are recognized: 1) Sensible-heat polynyas, caused by the upwelling of warmer water and 2) Latent-heat polynyas, resulting from persistent winds from the coastline.
Formation
Only the top layer of water needs to cool to the freezing point. Convection of the surface layer involves the top , down to the pycnocline of increased density.
In calm water, the first sea ice to form on the surface is a skim of separate crystals which initially are in the form of tiny discs, floating flat on the surface and of diameter less than . Each disc has its c-axis vertical and grows outwards laterally. At a certain point such a disc shape becomes unstable and the growing isolated crystals take on a hexagonal, stellar form, with long fragile arms stretching out over the surface. These crystals also have their c-axis vertical. The dendritic arms are very fragile and soon break off, leaving a mixture of discs and arm fragments. With any kind of turbulence in the water, these fragments break up further into random-shaped small crystals which form a suspension of increasing density in the surface water, an ice type called frazil or grease ice. In quiet conditions the frazil crystals soon freeze together to form a continuous thin sheet of young ice; in its early stages, when it is still transparent – that is the ice called nilas. Once nilas has formed, a quite different growth process occurs, in which water freezes on to the bottom of the existing ice sheet, a process called congelation growth. This growth process yields first-year ice.
In rough water, fresh sea ice is formed by the cooling of the ocean as heat is lost into the atmosphere. The uppermost layer of the ocean is supercooled to slightly below the freezing point, at which time tiny ice platelets (frazil ice) form. With time, this process leads to a mushy surface layer, known as grease ice. Frazil ice formation may also be started by snowfall, rather than supercooling. Waves and wind then act to compress these ice particles into larger plates, of several meters in diameter, called pancake ice. These float on the ocean surface and collide with one another, forming upturned edges. In time, the pancake ice plates may themselves be rafted over one another or frozen together into a more solid ice cover, known as consolidated pancake ice. Such ice has a very rough appearance on top and bottom.
If sufficient snow falls on sea ice to depress the freeboard below sea level, sea water will flow in and a layer of ice will form of mixed snow/sea water. This is particularly common around Antarctica.
Russian scientist Vladimir Vize (1886–1954) devoted his life to study the Arctic ice pack and developed the Scientific Prediction of Ice Conditions Theory, for which he was widely acclaimed in academic circles. He applied this theory in the field in the Kara Sea, which led to the discovery of Vize Island.
Yearly freeze and melt cycle
The annual freeze and melt cycle is set by the annual cycle of solar insolation and of ocean and atmospheric temperature and of variability in this annual cycle.
In the Arctic, the area of ocean covered by sea ice increases over winter from a minimum in September to a maximum in March or sometimes February, before melting over the summer. In the Antarctic, where the seasons are reversed, the annual minimum is typically in February and the annual maximum in September or October. The presence of sea ice abutting the calving fronts of ice shelves has been shown to influence glacier flow and potentially the stability of the Antarctic ice sheet.
The growth and melt rate are also affected by the state of the ice itself. During growth, the ice thickening due to freezing (as opposed to dynamics) is itself dependent on the thickness, so that the ice growth slows as the ice thickens. Likewise, during melt, thinner sea ice melts faster. This leads to different behaviour between multiyear and first year ice. In addition, melt ponds on the ice surface during the melt season lower the albedo such that more solar radiation is absorbed, leading to a feedback where melt is accelerated. The presence of melt ponds is affected by the permeability of the sea ice (i.e. whether meltwater can drain) and the topography of the sea ice surface (i.e. the presence of natural basins for the melt ponds to form in). First year ice is flatter than multiyear ice due to the lack of dynamic ridging, so ponds tend to have greater area. They also have lower albedo since they are on thinner ice, which blocks less of the solar radiation from reaching the dark ocean below.
Physical properties
Sea ice is a composite material made up of pure ice, liquid brine, air, and salt. The volumetric fractions of these components—ice, brine, and air—determine the key physical properties of sea ice, including thermal conductivity, heat capacity, latent heat, density, elastic modulus, and mechanical strength. Brine volume fraction depends on sea-ice salinity and temperature, while sea-ice salinity mainly depends on ice age and thickness. During the ice growth period, its bulk brine volume is typically below 5%. Air volume fraction during ice growth period is typically around 1–2 %, but may substantially increase upon ice warming. Air volume of sea ice in can be as high as 15 % in summer and 4 % in autumn. Both brine and air volumes influence sea-ice density values, which are typically around 840–910 kg/m3 for first-year ice. Sea-ice density is a significant source of errors in sea-ice thickness retrieval using radar and laser satellite altimetry, resulting in uncertainties of 0.3–0.4 m.
Monitoring and observations
Changes in sea ice conditions are best demonstrated by the rate of melting over time. A composite record of Arctic ice demonstrates that the floes' retreat began around 1900, experiencing more rapid melting beginning within the past 50 years. Satellite study of sea ice began in 1979 and became a much more reliable measure of long-term changes in sea ice. In comparison to the extended record, the sea-ice extent in the polar region by September 2007 was only half the recorded mass that had been estimated to exist within the 1950–1970 period.
Arctic sea ice extent ice hit an all-time low in September 2012, when the ice was determined to cover only 24% of the Arctic Ocean, offsetting the previous low of 29% in 2007. Predictions of when the first "ice free" Arctic summer might occur vary.
Antarctic sea ice extent gradually increased in the period of satellite observations, which began in 1979, until a rapid decline in southern hemisphere spring of 2016.
Effects of climate change
Sea ice provides an ecosystem for various polar species, particularly the polar bear, whose environment is being threatened as global warming causes the ice to melt more as the Earth's temperature gets warmer. Furthermore, the sea ice itself functions to help keep polar climates cool, since the ice exists in expansive enough amounts to maintain a cold environment. At this, sea ice's relationship with global warming is cyclical; the ice helps to maintain cool climates, but as the global temperature increases, the ice melts and is less effective in keeping those climates cold. The bright, shiny surface (albedo) of the ice also serves a role in maintaining cooler polar temperatures by reflecting much of the sunlight that hits it back into space. As the sea ice melts, its surface area shrinks, diminishing the size of the reflective surface and therefore causing the earth to absorb more of the sun's heat. As the ice melts it lowers the albedo thus causing more heat to be absorbed by the Earth and further increase the amount of melting ice. Though the size of the ice floes is affected by the seasons, even a small change in global temperature can greatly affect the amount of sea ice and due to the shrinking reflective surface that keeps the ocean cool, this sparks a cycle of ice shrinking and temperatures warming. As a result, the polar regions are the most susceptible places to climate change on the planet.
Furthermore, sea ice affects the movement of ocean waters. In the freezing process, much of the salt in ocean water is squeezed out of the frozen crystal formations, though some remains frozen in the ice. This salt becomes trapped beneath the sea ice, creating a higher concentration of salt in the water beneath ice floes. This concentration of salt contributes to the salinated water's density and this cold, denser water sinks to the bottom of the ocean. This cold water moves along the ocean floor towards the equator, while warmer water on the ocean surface moves in the direction of the poles. This is referred to as "conveyor belt motion" and is a regularly occurring process.
Modelling
In order to gain a better understanding about the variability, numerical sea ice models are used to perform sensitivity studies. The two main ingredients are the ice dynamics and the thermodynamical properties (see Sea ice emissivity modelling, Sea ice growth processes and Sea ice thickness). There are many sea ice model computer codes available for doing this, including the CICE numerical suite.
Many global climate models (GCMs) have sea ice implemented in their numerical simulation scheme in order to capture the ice–albedo feedback correctly. Examples include:
The Louvain-la-Neuve Sea Ice Model is a numerical model of sea ice designed for climate studies and operational oceanography developed at Université catholique de Louvain. It is coupled to the ocean general circulation model OPA (Ocean Parallélisé) and is freely available as a part of the Nucleus for European Modeling of the Ocean.
The MIT General Circulation Model is a global circulation model developed at Massachusetts Institute of Technology includes a package for sea-ice. The code is freely available there.
The University Corporation for Atmospheric Research develops the Community Sea Ice Model.
CICE is run by the Los Alamos National Laboratory. The project is open source and designed as a component of GCM, although it provides a standalone mode.
The Finite-Element Sea-Ice Ocean Model developed at Alfred Wegener Institute uses an unstructured grid.
The neXt Generation Sea-Ice model (neXtSIM) is a Lagrangian model using an adaptive and unstructured triangular mesh and includes a new and unique class of rheological model called Maxwell-Elasto-Brittle to treat the ice dynamics. This model is developed at the Nansen Center in Bergen, Norway.
The Coupled Model Intercomparison Project offers a standard protocol for studying the output of coupled atmosphere-ocean general circulation models. The coupling takes place at the atmosphere-ocean interface where the sea ice may occur.
In addition to global modeling, various regional models deal with sea ice. Regional models are employed for seasonal forecasting experiments and for process studies.
Ecology
Sea ice is part of the Earth's biosphere. When sea water freezes, the ice is riddled with brine-filled channels which sustain sympagic organisms such as bacteria, algae, copepods and annelids, which in turn provide food for animals such as krill and specialised fish like the bald notothen, fed upon in turn by larger animals such as emperor penguins and minke whales.
A decline of seasonal sea ice puts the survival of Arctic species such as ringed seals and polar bears at risk.
Extraterrestrial presence
Other element and compounds have been speculated to exist as oceans and seas on extraterrestrial planets. Scientists notably suspect the existence of "icebergs" of solid diamond and corresponding seas of liquid carbon on the ice giants, Neptune and Uranus. This is due to extreme pressure and heat at the core, that would turn carbon into a supercritical fluid.
| Physical sciences | Glaciology | null |
343230 | https://en.wikipedia.org/wiki/Autopilot | Autopilot | An autopilot is a system used to control the path of a vehicle without requiring constant manual control by a human operator. Autopilots do not replace human operators. Instead, the autopilot assists the operator's control of the vehicle, allowing the operator to focus on broader aspects of operations (for example, monitoring the trajectory, weather and on-board systems).
When present, an autopilot is often used in conjunction with an autothrottle, a system for controlling the power delivered by the engines.
An autopilot system is sometimes colloquially referred to as "George" (e.g. "we'll let George fly for a while"; "George is flying the plane now".). The etymology of the nickname is unclear: some claim it is a reference to American inventor George De Beeson (1897 - 1965), who patented an autopilot in the 1930s, while others claim that Royal Air Force pilots coined the term during World War II to symbolize that their aircraft technically belonged to King George VI.
First autopilots
In the early days of aviation, aircraft required the continuous attention of a pilot to fly safely. As aircraft range increased, allowing flights of many hours, the constant attention led to serious fatigue. An autopilot is designed to perform some of the pilot's tasks.
The first aircraft autopilot was developed by Sperry Corporation in 1912. The autopilot connected a gyroscopic heading indicator, and attitude indicator to hydraulically operated elevators and rudder. (Ailerons were not connected as wing dihedral was counted upon to produce the necessary roll stability.) It permitted the aircraft to fly straight and level on a compass course without a pilot's attention, greatly reducing the pilot's workload.
Lawrence Sperry, the son of famous inventor Elmer Sperry, demonstrated it in 1914 at an aviation safety contest held in Paris. Sperry demonstrated the credibility of the invention by flying the aircraft with his hands away from the controls and visible to onlookers. Elmer Sperry Jr., the son of Lawrence Sperry, and Capt Shiras continued work on the same autopilot after the war, and in 1930, they tested a more compact and reliable autopilot which kept a U.S. Army Air Corps aircraft on a true heading and altitude for three hours.
In 1930, the Royal Aircraft Establishment in the United Kingdom developed an autopilot called a pilots' assister that used a pneumatically spun gyroscope to move the flight controls.
The autopilot was further developed, to include, for example, improved control algorithms and hydraulic servomechanisms. Adding more instruments, such as radio-navigation aids, made it possible to fly at night and in bad weather. In 1947, a U.S. Air Force C-53 made a transatlantic flight, including takeoff and landing, completely under the control of an autopilot.
Bill Lear developed his F-5 automatic pilot, and automatic approach control system, and was awarded the Collier Trophy in 1949.
In the early 1920s, the Standard Oil tanker J.A. Moffet became the first ship to use an autopilot.
The Piasecki HUP-2 Retriever was the first production helicopter with an autopilot.
The lunar module digital autopilot of the Apollo program is an early example of a fully digital autopilot system in spacecraft.
Modern autopilots
Not all of the passenger aircraft flying today have an autopilot system. Older and smaller general aviation aircraft especially are still hand-flown, and even small airliners with fewer than twenty seats may also be without an autopilot as they are used on short-duration flights with two pilots. The installation of autopilots in aircraft with more than twenty seats is generally made mandatory by international aviation regulations. There are three levels of control in autopilots for smaller aircraft. A single-axis autopilot controls an aircraft in the roll axis only; such autopilots are also known colloquially as "wing levellers", reflecting their single capability. A two-axis autopilot controls an aircraft in the pitch axis as well as roll, and may be little more than a wing leveller with limited pitch oscillation-correcting ability; or it may receive inputs from on-board radio navigation systems to provide true automatic flight guidance once the aircraft has taken off until shortly before landing; or its capabilities may lie somewhere between these two extremes. A three-axis autopilot adds control in the yaw axis and is not required in many small aircraft.
Autopilots in modern complex aircraft are three-axis and generally divide a flight into taxi, takeoff, climb, cruise (level flight), descent, approach, and landing phases. Autopilots that automate all of these flight phases except taxi and takeoff exist. An autopilot-controlled approach to landing on a runway and controlling the aircraft on rollout (i.e. keeping it on the centre of the runway) is known as an Autoland, where the autopilot utilizes an Instrument Landing System (ILS) Cat IIIc approach, which is used when the visibility is zero. These approaches are available at many major airports' runways today, especially at airports subject to adverse weather phenomena such as fog. The aircraft can typically stop on their own, but will require the disengagement of the autopilot in order to exit the runway and taxi to the gate. An autopilot is often an integral component of a Flight Management System.
Modern autopilots use computer software to control the aircraft. The software reads the aircraft's current position, and then controls a flight control system to guide the aircraft. In such a system, besides classic flight controls, many autopilots incorporate thrust control capabilities that can control throttles to optimize the airspeed.
The autopilot in a modern large aircraft typically reads its position and the aircraft's attitude from an inertial guidance system. Inertial guidance systems accumulate errors over time. They will incorporate error reduction systems such as the carousel system that rotates once a minute so that any errors are dissipated in different directions and have an overall nulling effect. Error in gyroscopes is known as drift. This is due to physical properties within the system, be it mechanical or laser guided, that corrupt positional data. The disagreements between the two are resolved with digital signal processing, most often a six-dimensional Kalman filter. The six dimensions are usually roll, pitch, yaw, altitude, latitude, and longitude. Aircraft may fly routes that have a required performance factor, therefore the amount of error or actual performance factor must be monitored in order to fly those particular routes. The longer the flight, the more error accumulates within the system. Radio aids such as DME, DME updates, and GPS may be used to correct the aircraft position.
Control Wheel Steering
An option midway between fully automated flight and manual flying is Control Wheel Steering (CWS). Although it is becoming less used as a stand-alone option in modern airliners, CWS is still a function on many aircraft today. Generally, an autopilot that is CWS equipped has three positions: off, CWS, and CMD. In CMD (Command) mode the autopilot has full control of the aircraft, and receives its input from either the heading/altitude setting, radio and navaids, or the FMS (Flight Management System). In CWS mode, the pilot controls the autopilot through inputs on the yoke or the stick. These inputs are translated to a specific heading and attitude, which the autopilot will then hold until instructed to do otherwise. This provides stability in pitch and roll. Some aircraft employ a form of CWS even in manual mode, such as the MD-11 which uses a constant CWS in roll. In many ways, a modern Airbus fly-by-wire aircraft in Normal Law is always in CWS mode. The major difference is that in this system the limitations of the aircraft are guarded by the flight control computer, and the pilot cannot steer the aircraft past these limits.
Computer system details
The hardware of an autopilot varies between implementations, but is generally designed with redundancy and reliability as foremost considerations. For example, the Rockwell Collins AFDS-770 Autopilot Flight Director System used on the Boeing 777 uses triplicated FCP-2002 microprocessors which have been formally verified and are fabricated in a radiation-resistant process.
Software and hardware in an autopilot are tightly controlled, and extensive test procedures are put in place.
Some autopilots also use design diversity. In this safety feature, critical software processes will not only run on separate computers, and possibly even using different architectures, but each computer will run software created by different engineering teams, often being programmed in different programming languages. It is generally considered unlikely that different engineering teams will make the same mistakes. As the software becomes more expensive and complex, design diversity is becoming less common because fewer engineering companies can afford it. The flight control computers on the Space Shuttle used this design: there were five computers, four of which redundantly ran identical software, and a fifth backup running software that was developed independently. The software on the fifth system provided only the basic functions needed to fly the Shuttle, further reducing any possible commonality with the software running on the four primary systems.
Stability augmentation systems
A stability augmentation system (SAS) is another type of automatic flight control system; however, instead of maintaining the aircraft required altitude or flight path, the SAS will move the aircraft control surfaces to damp unacceptable motions. SAS automatically stabilizes the aircraft in one or more axes. The most common type of SAS is the yaw damper which is used to reduce the Dutch roll tendency of swept-wing aircraft. Some yaw dampers are part of the autopilot system while others are stand-alone systems.
Yaw dampers use a sensor to detect how fast the aircraft is rotating (either a gyroscope or a pair of accelerometers), a computer/amplifier and an actuator. The sensor detects when the aircraft begins the yawing part of Dutch roll. A computer processes the signal from the sensor to determine the rudder deflection required to damp the motion. The computer tells the actuator to move the rudder in the opposite direction to the motion since the rudder has to oppose the motion to reduce it. The Dutch roll is damped and the aircraft becomes stable about the yaw axis. Because Dutch roll is an instability that is inherent in all swept-wing aircraft, most swept-wing aircraft need some sort of yaw damper.
There are two types of yaw damper: the series yaw damper and the parallel yaw damper. The actuator of a parallel yaw damper will move the rudder independently of the pilot's rudder pedals while the actuator of a series yaw damper is clutched to the rudder control quadrant, and will result in pedal movement when the rudder moves.
Some aircraft have stability augmentation systems that will stabilize the aircraft in more than a single axis. The Boeing B-52, for example, requires both pitch and yaw SAS in order to provide a stable bombing platform. Many helicopters have pitch, roll and yaw SAS systems. Pitch and roll SAS systems operate much the same way as the yaw damper described above; however, instead of damping Dutch roll, they will damp pitch and roll oscillations to improve the overall stability of the aircraft.
Autopilot for ILS landings
Instrument-aided landings are defined in categories by the International Civil Aviation Organization, or ICAO. These are dependent upon the required visibility level and the degree to which the landing can be conducted automatically without input by the pilot.
CAT I – This category permits pilots to land with a decision height of and a forward visibility or Runway Visual Range (RVR) of . Autopilots are not required.
CAT II – This category permits pilots to land with a decision height between and and a RVR of . Autopilots have a fail passive requirement.
CAT IIIa -This category permits pilots to land with a decision height as low as and a RVR of . It needs a fail-passive autopilot. There must be only a 10−6 probability of landing outside the prescribed area.
CAT IIIb – As IIIa but with the addition of automatic roll out after touchdown incorporated with the pilot taking control some distance along the runway. This category permits pilots to land with a decision height less than 50 feet or no decision height and a forward visibility of in Europe (76 metres, compare this to aircraft size, some of which are now over long) or in the United States. For a landing-without-decision aid, a fail-operational autopilot is needed. For this category some form of runway guidance system is needed: at least fail-passive but it needs to be fail-operational for landing without decision height or for RVR below .
CAT IIIc – As IIIb but without decision height or visibility minimums, also known as "zero-zero". Not yet implemented as it would require the pilots to taxi in zero-zero visibility. An aircraft that is capable of landing in a CAT IIIb that is equipped with autobrake would be able to fully stop on the runway but would have no ability to taxi.
Fail-passive autopilot: in case of failure, the aircraft stays in a controllable position and the pilot can take control of it to go around or finish landing. It is usually a dual-channel system.
Fail-operational autopilot: in case of a failure below alert height, the approach, flare and landing can still be completed automatically. It is usually a triple-channel system or dual-dual system.
Radio-controlled models
In radio-controlled modelling, and especially RC aircraft and helicopters, an autopilot is usually a set of extra hardware and software that deals with pre-programming the model's flight.
Flight Director
A flight director (FD) is a flight instrument that is overlaid on the attitude indicator that shows the pilot of an aircraft the attitude required to execute the desired flight path. While the flight director is separate from the autopilot, they are closely linked. With a flight plan programmed into the flight computer, the flight director will command rolls when turns are required.
Without a flight director, the autopilot is limited to more basic modes, such as maintaining an altitude or a heading, or turning on to a new heading when commanded by the pilot.
When the autopilot and flight director are used together, more complex autopilot modes are possible. The autopilot can follow flight director commands, thus following the flight plan route without pilot intervention.
| Technology | Aircraft components | null |
343246 | https://en.wikipedia.org/wiki/Ice%20shelf | Ice shelf | An ice shelf is a large platform of glacial ice floating on the ocean, fed by one or multiple tributary glaciers. Ice shelves form along coastlines where the ice thickness is insufficient to displace the more dense surrounding ocean water. The boundary between the ice shelf (floating) and grounded ice (resting on bedrock or sediment) is referred to as the grounding line; the boundary between the ice shelf and the open ocean (often covered by sea ice) is the ice front or calving front.
Ice shelves are found in Antarctica and the Arctic (Greenland, Northern Canada, and the Russian Arctic), and can range in thickness from about . The world's largest ice shelves are the Ross Ice Shelf and the Filchner-Ronne Ice Shelf in Antarctica.
The movement of ice shelves is principally driven by gravity-induced pressure from the grounded ice. That flow continually moves ice from the grounding line to the seaward front of the shelf. Typically, a shelf front will extend forward for years or decades between major calving events (calving is the sudden release and breaking away of a mass of ice from a glacier, iceberg, ice front, ice shelf, or crevasse). Snow accumulation on the upper surface and melting from the lower surface are also important to the mass balance of an ice shelf. Ice may also accrete onto the underside of the shelf.
The effects of climate change are visible in the changes to the cryosphere, such as reduction in sea ice and ice sheets, and disruption of ice shelves. In the last several decades, glaciologists have observed consistent decreases in ice shelf extent through melt, calving, and complete disintegration of some shelves. Well studied examples include disruptions of the Thwaites Ice Shelf, Larsen Ice Shelf, Filchner–Ronne Ice Shelf (all three in the Antarctic) and the disruption of the Ellesmere Ice Shelf in the Arctic.
Definition
An ice shelf is "a floating slab of ice originating from land of considerable thickness extending from the coast (usually of great horizontal extent with a very gently sloping surface), resulting from the flow of ice sheets, initially formed by the accumulation of snow, and often filling embayments in the coastline of an ice sheet."
In contrast, sea ice is formed on water, is much thinner (typically less than ), and forms throughout the Arctic Ocean. It is also found in the Southern Ocean around the continent of Antarctica.
The term captured ice shelf has been used for the ice over a subglacial lake, such as Lake Vostok.
Properties
Ice shelves are thick plates of ice, formed continuously by glaciers, that float atop an ocean. The shelves act as "brakes" for the glaciers. These shelves serve another important purpose—"they moderate the amount of melting that occurs on the glaciers' surfaces. Once their ice shelves are removed, the glaciers increase in speed due to meltwater percolation and/or a reduction of braking forces, and they may begin to dump more ice into the ocean than they gather as snow in their catchments. Glacier ice speed increases are already observed in Peninsula areas where ice shelves disintegrated in prior years."
Height
The density contrast between glacial ice and liquid water means that at least of the floating ice is above the ocean surface, depending on how much pressurized air is contained in the bubbles within the glacial ice, stemming from compressed snow. The formula for the denominators above is , density of cold seawater is about 1028 kg/m3 and that of glacial ice from about 850 kg/m3 to well below 920 kg/m3, the limit for very cold ice without bubbles. The height of the shelf above the sea can be even larger, if there is much less dense firn and snow above the glacier ice.
By country or region
Antarctica
A large portion of the Antarctic coastline has ice shelves attached. Their aggregate area is over .
It has been found that of all the ice shelves on Earth, nearly all of them are in Antarctica.
In steady state, about half of Antarctica's ice shelf mass is lost to basal melt and half is lost to calving, but the relative importance of each process varies significantly between ice shelves. In recent decades, Antarctica's ice shelves have been out of balance, as they have lost more mass to basal melt and calving than has been replenished by the influx of new ice and snow.
Ross Ice Shelf
Filchner–Ronne Ice Shelf
Arctic
Canada
All Canadian ice shelves are attached to Ellesmere Island and lie north of 82°N. Ice shelves that are still in existence are the Alfred Ernest Ice Shelf, Ward Hunt Ice Shelf, Milne Ice Shelf and Smith Ice Shelf. The M'Clintock Ice Shelf broke up from 1963 to 1966; the Ayles Ice Shelf broke up in 2005; and the Markham Ice Shelf broke up in 2008. The remaining ice shelves have also lost a significant amount of their area over time, with the Milne Ice Shelf being the last to be affected, with it breaking off in August 2020.
Russia
The Matusevich Ice Shelf was a ice shelf located in Severnaya Zemlya being fed by some of the largest ice caps on October Revolution Island, the Karpinsky Ice Cap to the south and the Rusanov Ice Cap to the north. In 2012 it ceased to exist.
Disruption due to climate change
In the last several decades, glaciologists have observed consistent decreases in ice shelf extent through melt, calving, and complete disintegration of some shelves. Well studied examples include disruptions of the Thwaites Ice Shelf, Larsen Ice Shelf, Filchner–Ronne Ice Shelf (all three in the Antarctic) and the disruption of the Ellesmere Ice Shelf in the Arctic.
The effects of climate change are visible in the changes to the cryosphere, such as reduction in sea ice and ice sheets, and disruption of ice shelves.
Disruption of Thwaites Ice Shelf
Disruption of Larsen Ice Shelf
Two sections of Antarctica's Larsen Ice Shelf broke apart into hundreds of unusually small fragments (hundreds of meters wide or less) in 1995 and 2002, Larsen C calved a huge ice island in 2017.
Disruption of Larsen B Ice Shelf
Disruption of Filchner–Ronne Ice Shelf
Other ice shelves in Antarctica
Wordie Ice Shelf has gone from an area of in 1950 to in 2000.
Prince Gustav Ice Shelf has gone from an area of to in 2008. After their loss the reduced buttressing of feeder glaciers has allowed the expected speed-up of inland ice masses after shelf ice break-up.
The Ross Ice Shelf is the largest ice shelf of Antarctica (an area of roughly and about across: about the size of France).
Wilkins Ice Shelf is another ice shelf that has suffered substantial retreat. The ice shelf had an area of in 1998 when was lost that year. In 2007 and 2008 significant rifting developed and led to the loss of another of area and some of the calving occurred in the Austral winter. The calving seemed to have resulted from preconditioning such as thinning, possibly due to basal melt, as surface melt was not as evident, leading to a reduction in the strength of the pinning point connections. The thinner ice then experienced spreading rifts and breakup. This period culminated in the collapse of an ice bridge connecting the main ice shelf to Charcot Island leading to the loss of an additional between February and June 2009.
Disruption of Ellesmere Ice Shelf (Arctic)
The Ellesmere ice shelf was reduced by 90% in the twentieth century, leaving the separate Alfred Ernest, Ayles, Milne, Ward Hunt, and Markham ice shelves. A 1986 survey of Canadian ice shelves found that 48 km2 (3.3 cubic kilometres) of ice calved from the Milne and Ayles ice shelves between 1959 and 1974. The Ayles Ice Shelf calved entirely on August 13, 2005. The Ward Hunt Ice Shelf, the largest remaining section of thick (>) landfast sea ice along the northern coastline of Ellesmere Island, lost of ice in a massive calving in 1961–1962. It further decreased by 27% in thickness () between 1967 and 1999. In the summer of 2002, the Ward Ice Shelf experienced another major breakup, and other instances of note happened in 2008 and 2010 as well. The last remnant to remain mostly intact, the Milne Ice Shelf, also ultimately experienced a major breakup at the end of July 2020, losing over 40% of its area.
| Physical sciences | Glaciology | null |
343492 | https://en.wikipedia.org/wiki/Electrolytic%20capacitor | Electrolytic capacitor | An electrolytic capacitor is a polarized capacitor whose anode or positive plate is made of a metal that forms an insulating oxide layer through anodization. This oxide layer acts as the dielectric of the capacitor. A solid, liquid, or gel electrolyte covers the surface of this oxide layer, serving as the cathode or negative plate of the capacitor. Because of their very thin dielectric oxide layer and enlarged anode surface, electrolytic capacitors have a much higher capacitance-voltage (CV) product per unit volume than ceramic capacitors or film capacitors, and so can have large capacitance values. There are three families of electrolytic capacitor: aluminium electrolytic capacitors, tantalum electrolytic capacitors, and niobium electrolytic capacitors.
The large capacitance of electrolytic capacitors makes them particularly suitable for passing or bypassing low-frequency signals, and for storing large amounts of energy. They are widely used for decoupling or noise filtering in power supplies and DC link circuits for variable-frequency drives, for coupling signals between amplifier stages, and storing energy as in a flashlamp.
Electrolytic capacitors are polarized components because of their asymmetrical construction and must be operated with a higher potential (i.e., more positive) on the anode than on the cathode at all times. For this reason the polarity is marked on the device housing. Applying a reverse polarity voltage, or a voltage exceeding the maximum rated working voltage of as little as 1 or 1.5 volts, can damage the dielectric causing catastrophic failure of the capacitor itself. Failure of electrolytic capacitors can result in an explosion or fire, potentially causing damage to other components as well as injuries. Bipolar electrolytic capacitors which may be operated with either polarity are also made, using special constructions with two anodes connected in series. A bipolar electrolytic capacitor can be made by connecting two normal electrolytic capacitors in series, anode to anode or cathode to cathode, along with diodes.
General information
Electrolytic capacitors family tree
As to the basic construction principles of electrolytic capacitors, there are three different types: aluminium, tantalum, and niobium capacitors. Each of these three capacitor families uses non-solid and solid manganese dioxide or solid polymer electrolytes, so a great spread of different combinations of anode material and solid or non-solid electrolytes is available.
Charge principle
Like other conventional capacitors, electrolytic capacitors store the electric energy statically by charge separation in an electric field in the dielectric oxide layer between two electrodes. The non-solid or solid electrolyte in principle is the cathode, which thus forms the second electrode of the capacitor. This and the storage principle distinguish them from electrochemical capacitors or supercapacitors, in which the electrolyte generally is the ionic conductive connection between two electrodes and the storage occurs with statically double-layer capacitance and electrochemical pseudocapacitance.
Basic materials and construction
Electrolytic capacitors use a chemical feature of some special metals, previously called "valve metals", which on contact with a particular electrolyte form a very thin insulating oxide layer on their surface by anodic oxidation which can function as a dielectric. There are three different anode metals in use for electrolytic capacitors:
Aluminum electrolytic capacitors use a high-purity etched aluminium foil with aluminium oxide as dielectric
Tantalum electrolytic capacitors use a sintered pellet (“slug”) of high-purity tantalum powder with tantalum pentoxide as dielectric
Niobium electrolytic capacitors use a sintered "slug" of high-purity niobium or niobium oxide powder with niobium pentoxide as dielectric.
To increase their capacitance per unit volume, all anode materials are either etched or sintered and have a rough surface structure with a much higher surface area compared to a smooth surface of the same area or the same volume. By applying a positive voltage to the above-mentioned anode material in an electrolytic bath an oxide barrier layer with a thickness corresponding to the applied voltage will be formed (formation). This oxide layer acts as the dielectric in an electrolytic capacitor. The properties of these oxide layers are given in the following table:
After forming a dielectric oxide on the rough anode structure, a counter electrode has to match the rough insulating oxide surface. This is accomplished by the electrolyte, which acts as the cathode electrode of an electrolytic capacitor. There are many different electrolytes in use. Generally they are distinguished into two species, “non-solid” and “solid” electrolytes. As a liquid medium which has ion conductivity caused by moving ions, non-solid electrolytes can easily fit the rough structures. Solid electrolytes which have electron conductivity can fit the rough structures with the help of special chemical processes like pyrolysis for manganese dioxide or polymerization for conducting polymers.
Comparing the permittivities of the different oxide materials it is seen that tantalum pentoxide has a permittivity approximately three times higher than aluminium oxide. Tantalum electrolytic capacitors of a given CV value theoretically are therefore smaller than aluminium electrolytic capacitors. In practice different safety margins to reach reliable components makes a comparison difficult.
The anodically generated insulating oxide layer is destroyed if the polarity of the applied voltage changes.
Capacitance and volumetric efficiency
Electrolytic capacitors are based on the principle of a "plate capacitor" whose capacitance increases with larger electrode area A, higher dielectric permittivity ε, and thinness of dielectric (d).
The dielectric thickness of electrolytic capacitors is very small, in the range of nanometers per volt. On the other hand, the voltage strengths of these oxide layers are quite high. With this very thin dielectric oxide layer combined with a sufficiently high dielectric strength the electrolytic capacitors can achieve a high volumetric capacitance. This is one reason for the high capacitance values of electrolytic capacitors compared to conventional capacitors.
All etched or sintered anodes have a much higher surface area compared to a smooth surface of the same area or the same volume. That increases the capacitance value, depending on the rated voltage, by a factor of up to 200 for non-solid aluminium electrolytic capacitors as well as for solid tantalum electrolytic capacitors. The large surface compared to a smooth one is the second reason for the relatively high capacitance values of electrolytic capacitors compared with other capacitor families.
Because the forming voltage defines the oxide layer thickness, the desired voltage rating can be produced very simply. Electrolytic capacitors have high volumetric efficiency, the so-called "CV product", defined as the product of capacitance and voltage divided by volume.
Basic construction of non-solid aluminium electrolytic capacitors
Basic construction of solid tantalum electrolytic capacitors
Types and features of electrolytic capacitors
Comparison of electrolytic capacitor types
Combinations of anode materials for electrolytic capacitors and the electrolytes used have given rise to wide varieties of capacitor types with different properties. An outline of the main characteristics of the different types is shown in the table below.
The non-solid or so-called "wet" aluminium electrolytic capacitors were and are the cheapest among all other conventional capacitors. They not only provide the cheapest solutions for high capacitance or voltage values for decoupling and buffering purposes but are also insensitive to low ohmic charging and discharging as well as to low-energy transients. Non-solid electrolytic capacitors can be found in nearly all areas of electronic devices, with the exception of military applications.
Tantalum electrolytic capacitors with solid electrolyte as surface-mountable chip capacitors are mainly used in electronic devices in which little space is available or a low profile is required. They operate reliably over a wide temperature range without large parameter deviations. In military and space applications only tantalum electrolytic capacitors have the necessary approvals.
Niobium electrolytic capacitors are in direct competition with industrial tantalum electrolytic capacitors because niobium is more readily available. Their properties are comparable.
The electrical properties of aluminium, tantalum and niobium electrolytic capacitors have been greatly improved by the polymer electrolyte.
Comparison of electrical parameters
In order to compare the different characteristics of the different electrolytic capacitor types, capacitors with the same dimensions and of similar capacitance and voltage are compared in the following table. In such a comparison the values for ESR and ripple current load are the most important parameters for the use of electrolytic capacitors in modern electronic equipment. The lower the ESR, the higher the ripple current per volume and better functionality of the capacitor in the circuit. However, better electrical parameters come with higher prices.
1) Manufacturer, series name, capacitance/voltage
2) calculated for a capacitor 100 μF/10 V,
3) from a 1976 data sheet
Styles of aluminium and tantalum electrolytic capacitors
Aluminium electrolytic capacitors form the bulk of the electrolytic capacitors used in electronics because of the large diversity of sizes and the inexpensive production. Tantalum electrolytic capacitors, usually used in the SMD (surface-mount device) version, have a higher specific capacitance than the aluminium electrolytic capacitors and are used in devices with limited space or flat design such as laptops. They are also used in military technology, mostly in axial style, hermetically sealed. Niobium electrolytic chip capacitors are a new development in the market and are intended as a replacement for tantalum electrolytic chip capacitors.
History
Origin
The phenomenon that in an electrochemical process, aluminium and such metals as tantalum, niobium, manganese, titanium, zinc, cadmium, etc., can form an oxide layer which blocks an electric current from flowing in one direction but which allows current to flow in the opposite direction, was first observed in 1857 by the German physicist and chemist Johann Heinrich Buff (1805–1878). It was first put to use in 1875 by the French researcher and founder Eugène Ducretet, who coined the term "valve metal" for such metals.
Charles Pollak (born Karol Pollak), a producer of accumulators, found out that the oxide layer on an aluminium anode remained stable in a neutral or alkaline electrolyte, even when the power was switched off. In 1896, he filed a patent for an "Electric liquid capacitor with aluminium electrodes" (de: Elektrischer Flüssigkeitskondensator mit Aluminiumelektroden) based on his idea of using the oxide layer in a polarized capacitor in combination with a neutral or slightly alkaline electrolyte.
"Wet" aluminium capacitor
The first industrially realized electrolytic capacitors consisted of a metallic box used as the cathode. It was filled with a borax electrolyte dissolved in water, in which a folded aluminium anode plate was inserted. Applying a DC voltage from outside, an oxide layer was formed on the surface of the anode. The advantage of these capacitors was that they were significantly smaller and cheaper than all other capacitors at this time relative to the realized capacitance value. This construction with different styles of anode construction but with a case as cathode and container for the electrolyte was used up to the 1930s and was called a "wet" electrolytic capacitor, in the sense of its having a high water content.
The first more common application of wet aluminium electrolytic capacitors was in large telephone exchanges, to reduce relay hash (noise) on the 48 volt DC power supply. The development of AC-operated domestic radio receivers in the late 1920s created a demand for large-capacitance (for the time) and high-voltage capacitors for the valve amplifier technique, typically at least 4 microfarads and rated at around 500 volts DC. Waxed paper and oiled silk film capacitors were available, but devices with that order of capacitance and voltage rating were bulky and prohibitively expensive.
"Dry" aluminium capacitor
The ancestor of the modern electrolytic capacitor was patented by Samuel Ruben in 1925, who teamed with Philip Mallory, the founder of the battery company that is now known as Duracell International. Ruben's idea adopted the stacked construction of a silver mica capacitor. He introduced a separated second foil to contact the electrolyte adjacent to the anode foil instead of using the electrolyte-filled container as the capacitor's cathode. The stacked second foil got its own terminal additional to the anode terminal and the container no longer had an electrical function. This type of electrolytic capacitor combined with a liquid or gel-like electrolyte of a non-aqueous nature, which is therefore dry in the sense of having a very low water content, became known as the "dry" type of electrolytic capacitor.
With Ruben's invention, together with the invention of wound foils separated with a paper spacer 1927 by A. Eckel of Hydra-Werke (Germany), the actual development of electrolytic capacitors began.
William Dubilier, whose first patent for electrolytic capacitors was filed in 1928, industrialized the new ideas for electrolytic capacitors and started the first large commercial production in 1931 in the Cornell-Dubilier (CD) factory in Plainfield, New Jersey. At the same time in Berlin, Germany, the "Hydra-Werke", an AEG company, started the production of electrolytic capacitors in large quantities. Another manufacturer, Ralph D. Mershon, had success in servicing the radio-market demand for electrolytic capacitors.
In his 1896 patent Pollak already recognized that the capacitance of the capacitor increases when roughening the surface of the anode foil. Today (2014), electrochemically etched low voltage foils can achieve an up to 200-fold increase in surface area compared to a smooth surface. Advances in the etching process are the reason for the dimension reductions in aluminium electrolytic capacitors over recent decades.
For aluminium electrolytic capacitors the decades from 1970 to 1990 were marked by the development of various new professional series specifically suited to certain industrial applications, for example with very low leakage currents or with long life characteristics, or for higher temperatures up to 125 °C.
Tantalum capacitors
One of the first tantalum electrolytic capacitors were developed in 1930 by Tansitor Electronic Inc. USA, for military purposes. The basic construction of a wound cell was adopted and a tantalum anode foil was used together with a tantalum cathode foil, separated with a paper spacer impregnated with a liquid electrolyte, mostly sulfuric acid, and encapsulated in a silver case.
The relevant development of solid electrolyte tantalum capacitors began some years after William Shockley, John Bardeen and Walter Houser Brattain invented the transistor in 1947. It was invented by Bell Laboratories in the early 1950s as a miniaturized, more reliable low-voltage support capacitor to complement their newly invented transistor. The solution found by R. L. Taylor and H. E. Haring at Bell Labs in early 1950 was based on experience with ceramics. They ground tantalum to a powder, which they pressed into a cylindrical form and then sintered at a high temperature between 1500 and 2000 °C under vacuum conditions, to produce a pellet ("slug").
These first sintered tantalum capacitors used a non-solid electrolyte, which does not fit the concept of solid electronics. In 1952 a targeted search at Bell Labs by D. A. McLean and F. S. Power for a solid electrolyte led to the invention of manganese dioxide as a solid electrolyte for a sintered tantalum capacitor.
Although fundamental inventions came from Bell Labs, the inventions for manufacturing commercially viable tantalum electrolytic capacitors came from researchers at the Sprague Electric Company. Preston Robinson, Sprague's Director of Research, is considered to be the actual inventor of tantalum capacitors in 1954. His invention was supported by R. J. Millard, who introduced the "reform" step in 1955, a significant improvement in which the dielectric of the capacitor was repaired after each dip-and-convert cycle of MnO2 deposition, which dramatically reduced the leakage current of the finished capacitors.
Although solid tantalum capacitors offered capacitors with lower ESR and leakage current values than the aluminium electrolytic capacitors, a 1980 price shock for tantalum dramatically reduced the applications of tantalum electrolytic capacitors, especially in the entertainment industry. The industry switched back to using aluminium electrolytic capacitors.
Solid electrolytes
The first solid electrolyte of manganese dioxide developed 1952 for tantalum capacitors had a conductivity 10 times better than all other types of non-solid electrolytes. It also influenced the development of aluminium electrolytic capacitors. In 1964 the first aluminium electrolytic capacitors with solid electrolyte SAL electrolytic capacitor came on the market, developed by Philips.
With the beginning of digitalization, Intel launched its first microcomputer, the MCS 4, in 1971. In 1972 Hewlett Packard launched one of the first pocket calculators, the HP 35. The requirements for capacitors increased in terms of lowering the equivalent series resistance (ESR) for bypass and decoupling capacitors.
It was not until 1983 when a new step toward ESR reduction was taken by Sanyo with its "OS-CON" aluminium electrolytic capacitors. These capacitors used a solid organic conductor, the charge transfer salt TTF-TCNQ (tetracyanoquinodimethane), which provided an improvement in conductivity by a factor of 10 compared with the manganese dioxide electrolyte.
The next step in ESR reduction was the development of conducting polymers by Alan J. Heeger, Alan MacDiarmid and Hideki Shirakawa in 1975. The conductivity of conductive polymers such as polypyrrole (PPy) or PEDOT is better than that of TCNQ by a factor of 100 to 500, and close to the conductivity of metals.
In 1991 Panasonic released its "SP-Cap", series of polymer aluminium electrolytic capacitors. These aluminium electrolytic capacitors with polymer electrolytes reached very low ESR values directly comparable to ceramic multilayer capacitors (MLCCs). They were still less expensive than tantalum capacitors and with their flat design for laptops and cell phones competed with tantalum chip capacitors as well.
Tantalum electrolytic capacitors with PPy polymer electrolyte cathode followed three years later. In 1993 NEC introduced its SMD polymer tantalum electrolytic capacitors, called "NeoCap". In 1997 Sanyo followed with the "POSCAP" polymer tantalum chips.
A new conductive polymer for tantalum polymer capacitors was presented by Kemet at the "1999 Carts" conference. This capacitor used the newly developed organic conductive polymer PEDT Poly(3,4-ethylenedioxythiophene), also known as PEDOT (trade name Baytron®)
Niobium capacitors
Another price explosion for tantalum in 2000/2001 forced the development of niobium electrolytic capacitors with manganese dioxide electrolyte, which have been available since 2002. Niobium is a sister metal to tantalum and serves as valve metal generating an oxide layer during anodic oxidation. Niobium as raw material is much more abundant in nature than tantalum and is less expensive. It was a question of the availability of the base metal in the late 1960s which led to development and implementation of niobium electrolytic capacitors in the former Soviet Union instead of tantalum capacitors as in the West. The materials and processes used to produce niobium-dielectric capacitors are essentially the same as for existing tantalum-dielectric capacitors. The characteristics of niobium electrolytic capacitors and tantalum electrolytic capacitors are roughly comparable.
Water-based electrolytes
With the goal of reducing ESR for inexpensive non-solid electrolytic capacitors from the mid-1980s in Japan, new water-based electrolytes for aluminium electrolytic capacitors were developed. Water is inexpensive, an effective solvent for electrolytes, and significantly improves the conductivity of the electrolyte. The Japanese manufacturer Rubycon was a leader in the development of new water-based electrolyte systems with enhanced conductivity in the late 1990s. The new series of non-solid electrolytic capacitors with water-based electrolyte was described in the data sheets as having "low ESR", "low impedance", "ultra-low impedance" or "high ripple current".
From 1999 through at least 2010, a stolen recipe for such a water-based electrolyte, in which important stabilizers were absent, led to the widespread problem of "bad caps" (failing electrolytic capacitors), leaking or occasionally bursting in computers, power supplies, and other electronic equipment, which became known as the "capacitor plague". In these electrolytic capacitors the water reacts quite aggressively with aluminium, accompanied by strong heat and gas development in the capacitor, resulting in premature equipment failure, and development of a cottage repair industry.
Electrical characteristics
Series-equivalent circuit
The electrical characteristics of capacitors are harmonized by the international generic specification IEC 60384-1. In this standard, the electrical characteristics of capacitors are described by an idealized series-equivalent circuit with electrical components which model all ohmic losses, capacitive and inductive parameters of an electrolytic capacitor:
C, the capacitance of the capacitor
RESR, the equivalent series resistance which summarizes all ohmic losses of the capacitor, usually abbreviated as "ESR"
LESL, the equivalent series inductance which is the effective self-inductance of the capacitor, usually abbreviated as "ESL".
Rleak, the resistance representing the leakage current of the capacitor
Capacitance, standard values and tolerances
The electrical characteristics of electrolytic capacitors depend on the structure of the anode and the electrolyte used. This influences the capacitance value of electrolytic capacitors, which depends on measuring frequency and temperature. Electrolytic capacitors with non-solid electrolytes show a broader aberration over frequency and temperature ranges than do capacitors with solid electrolytes.
The basic unit of an electrolytic capacitor's capacitance is the microfarad (μF). The capacitance value specified in the data sheets of the manufacturers is called the rated capacitance CR or nominal capacitance CN and is the value for which the capacitor has been designed.
The standardized measuring condition for electrolytic capacitors is an AC measuring method with 0.5 V at a frequency of 100/120 Hz at a temperature of 20 °C. For tantalum capacitors a DC bias voltage of 1.1 to 1.5 V for types with a rated voltage ≤2.5 V, or 2.1 to 2.5 V for types with a rated voltage of >2.5 V, may be applied during the measurement to avoid reverse voltage.
The capacitance value measured at the frequency of 1 kHz is about 10% less than the 100/120 Hz value. Therefore, the capacitance values of electrolytic capacitors are not directly comparable and differ from those of film capacitors or ceramic capacitors, whose capacitance is measured at 1 kHz or higher.
Measured with an AC measuring method at 100/120 Hz the capacitance value is the closest value to the electrical charge stored in the e-caps. The stored charge is measured with a special discharge method and is called the DC capacitance. The DC capacitance is about 10% higher than the 100/120 Hz AC capacitance. The DC capacitance is of interest for discharge applications like photoflash.
The percentage of allowed deviation of the measured capacitance from the rated value is called the capacitance tolerance. Electrolytic capacitors are available in different tolerance series, whose values are specified in the E series specified in IEC 60063. For abbreviated marking in tight spaces, a letter code for each tolerance is specified in IEC 60062.
rated capacitance, series E3, tolerance ±20%, letter code "M"
rated capacitance, series E6, tolerance ±20%, letter code "M"
rated capacitance, series E12, tolerance ±10%, letter code "K"
The required capacitance tolerance is determined by the particular application. Electrolytic capacitors, which are often used for filtering and bypassing, do not have the need for narrow tolerances because they are mostly not used for accurate frequency applications like in oscillators.
Rated and category voltage
Referring to the IEC/EN 60384-1 standard, the allowed operating voltage for electrolytic capacitors is called the "rated voltage UR" or "nominal voltage UN". The rated voltage UR is the maximum DC voltage or peak pulse voltage that may be applied continuously at any temperature within the rated temperature range TR.
The voltage proof of electrolytic capacitors decreases with increasing temperature. For some applications it is important to use a higher temperature range. Lowering the voltage applied at a higher temperature maintains safety margins. For some capacitor types therefore the IEC standard specifies a "temperature derated voltage" for a higher temperature, the "category voltage UC". The category voltage is the maximum DC voltage or peak pulse voltage that may be applied continuously to a capacitor at any temperature within the category temperature range TC. The relation between both voltages and temperatures is given in the picture at right.
Applying a higher voltage than specified may destroy electrolytic capacitors.
Applying a lower voltage may have a positive influence on electrolytic capacitors. For aluminium electrolytic capacitors a lower applied voltage can in some cases extend the lifetime. For tantalum electrolytic capacitors lowering the voltage applied increases the reliability and reduces the expected failure rate.
I
Surge voltage
The surge voltage indicates the maximum peak voltage value that may be applied to electrolytic capacitors during their application for a limited number of cycles.
The surge voltage is standardized in IEC/EN 60384-1. For aluminium electrolytic capacitors with a rated voltage of up to 315 V, the surge voltage is 1.15 times the rated voltage, and for capacitors with a rated voltage exceeding 315 V, the surge voltage is 1.10 times the rated voltage.
For tantalum electrolytic capacitors the surge voltage can be 1.3 times the rated voltage, rounded off to the nearest volt. The surge voltage applied to tantalum capacitors may influence the capacitor's failure rate.
Transient voltage
aluminium electrolytic capacitors with non-solid electrolyte are relatively insensitive to high and short-term transient voltages higher than surge voltage, if the frequency and the energy content of the transients are low. This ability depends on rated voltage and component size. Low energy transient voltages lead to a voltage limitation similar to a zener diode. An unambiguous and general specification of tolerable transients or peak voltages is not possible. In every case transients arise, the application has to be approved very carefully.
Electrolytic capacitors with solid manganese oxide or polymer electrolyte, and aluminium as well as tantalum electrolytic capacitors cannot withstand transients or peak voltages higher than the surge voltage. Transients may destroy this type of electrolytic capacitor.
Reverse voltage
Standard electrolytic capacitors, and aluminium as well as tantalum and niobium electrolytic capacitors are polarized and generally require the anode electrode voltage to be positive relative to the cathode voltage.
Nevertheless, electrolytic capacitors can withstand for short instants a reverse voltage for a limited number of cycles. Specifically, aluminium electrolytic capacitors with non-solid electrolyte can withstand a reverse voltage of about 1 V to 1.5 V. This reverse voltage should never be used to determine the maximum reverse voltage under which a capacitor can be used permanently.
Solid tantalum capacitors can also withstand reverse voltages for short periods. The most common guidelines for tantalum reverse voltage are:
10 % of rated voltage to a maximum of 1 V at 25 °C,
3 % of rated voltage to a maximum of 0.5 V at 85 °C,
1 % of rated voltage to a maximum of 0.1 V at 125 °C.
These guidelines apply for short excursion and should never be used to determine the maximum reverse voltage under which a capacitor can be used permanently.
But in no case, for aluminium as well as for tantalum and niobium electrolytic capacitors, may a reverse voltage be used for a permanent AC application.
To minimize the likelihood of a polarized electrolytic being incorrectly inserted into a circuit, polarity has to be very clearly indicated on the case, see the section on polarity marking below.
Special bipolar aluminium electrolytic capacitors designed for bipolar operation are available, and usually referred to as "non-polarized" or "bipolar" types. In these, the capacitors have two anode foils with full-thickness oxide layers connected in reverse polarity. On the alternate halves of the AC cycles, one of the oxides on the foil acts as a blocking dielectric, preventing reverse current from damaging the electrolyte of the other one. But these bipolar electrolytic capacitors are not suitable for main AC applications instead of power capacitors with metallized polymer film or paper dielectric.
Impedance
In general, a capacitor is seen as a storage component for electric energy. But this is only one capacitor application. A capacitor can also act as an AC resistor. aluminium electrolytic capacitors in particular are often used as decoupling capacitors to filter or bypass undesired AC frequencies to ground or for capacitive coupling of audio AC signals. Then the dielectric is used only for blocking DC. For such applications, the impedance (AC resistance) is as important as the capacitance value.
The impedance Z is the vector sum of reactance and resistance; it describes the phase difference and the ratio of amplitudes between sinusoidally varying voltage and sinusoidally varying current at a given frequency. In this sense impedance is a measure of the ability of the capacitor to pass alternating currents and can be used like Ohm's law.
In other words, impedance is a frequency-dependent AC resistance and possesses both magnitude and phase at a particular frequency.
In data sheets of electrolytic capacitors only the impedance magnitude |Z| is specified, and simply written as "Z". Regarding the IEC/EN 60384-1 standard, the impedance values of electrolytic capacitors are measured and specified at 10 kHz or 100 kHz depending on the capacitance and voltage of the capacitor.
Besides measuring, the impedance can be calculated using the idealized components of a capacitor's series-equivalent circuit, including an ideal capacitor C, a resistor ESR, and an inductance ESL. In this case the impedance at the angular frequency ω is given by the geometric (complex) addition of ESR, by a capacitive reactance XC
and by an inductive reactance XL (Inductance)
.
Then Z is given by
.
In the special case of resonance, in which the both reactive resistances XC and XL have the same value (XC=XL), then the impedance will only be determined by ESR. With frequencies above the resonance frequency, the impedance increases again because of the ESL of the capacitor. The capacitor becomes an inductor.
ESR and dissipation factor tan δ
The equivalent series resistance (ESR) summarizes all resistive losses of the capacitor. These are the terminal resistances, the contact resistance of the electrode contact, the line resistance of the electrodes, the electrolyte resistance, and the dielectric losses in the dielectric oxide layer.
For electrolytic capacitors, ESR generally decreases with increasing frequency and temperature.
ESR influences the superimposed AC ripple after smoothing and may influence the circuit functionality. Within the capacitor, ESR accounts for internal heat generation if a ripple current flows across the capacitor. This internal heat reduces the lifetime of non-solid aluminium electrolytic capacitors and affects the reliability of solid tantalum electrolytic capacitors.
For electrolytic capacitors, for historical reasons the dissipation factor tan δ will sometimes be specified in the data sheet instead of the ESR. The dissipation factor is determined by the tangent of the phase angle between the capacitive reactance XC minus the inductive reactance XL and the ESR. If the inductance ESL is small, the dissipation factor can be approximated as:
The dissipation factor is used for capacitors with very low losses in frequency-determining circuits where the reciprocal value of the dissipation factor is called the quality factor (Q), which represents a resonator's bandwidth.
Ripple current
"Ripple current" is the RMS value of a superimposed AC current of any frequency and any waveform of the current curve for continuous operation within the specified temperature range. It arises mainly in power supplies (including switched-mode power supplies) after rectifying an AC voltage and flows as charge and discharge current through any decoupling and smoothing capacitors.
Ripple currents generate heat inside the capacitor body. This dissipation power loss PL is caused by ESR and is the squared value of the effective (RMS) ripple current IR.
This internally generated heat, additional to the ambient temperature and possibly other external heat sources, leads to a capacitor body temperature having a temperature difference of Δ T relative to ambient. This heat has to be distributed as thermal losses Pth over the capacitor's surface A and the thermal resistance β to ambient.
The internally generated heat has to be distributed to ambient by thermal radiation, convection, and thermal conduction. The temperature of the capacitor, which is the net difference between heat produced and heat dissipated, must not exceed the capacitor's maximum specified temperature.
The ripple current is specified as an effective (RMS) value at 100 or 120 Hz or at 10 kHz at upper category temperature. Non-sinusoidal ripple currents have to be analyzed and separated into their single sinusoidal frequencies by means of Fourier analysis and summarized by squared addition the single currents.
In non-solid electrolytic capacitors the heat generated by the ripple current causes the evaporation of electrolytes, shortening the lifetime of the capacitors. Exceeding the limit tends to result in explosive failure.
In solid tantalum electrolytic capacitors with manganese dioxide electrolyte the heat generated by the ripple current affects the reliability of the capacitors. Exceeding the limit tends to result in catastrophic failure, failing short-circuit, with visible burning.
The heat generated by the ripple current also affects the lifetime of aluminium and tantalum electrolytic capacitors with solid polymer electrolytes. Exceeding the limit tends to result in catastrophic failure, failing short-circuit.
Current surge, peak or pulse current
aluminium electrolytic capacitors with non-solid electrolytes normally can be charged up to the rated voltage without any current surge, peak or pulse limitation. This property is a result of the limited ion movability in the liquid electrolyte, which slows down the voltage ramp across the dielectric, and of the capacitor's ESR. Only the frequency of peaks integrated over time must not exceed the maximal specified ripple current.
Solid tantalum electrolytic capacitors with manganese dioxide electrolyte or polymer electrolyte are damaged by peak or pulse currents. Solid Tantalum capacitors which are exposed to surge, peak or pulse currents, for example, in highly inductive circuits, should be used with a voltage derating. If possible, the voltage profile should be a ramp turn-on, as this reduces the peak current experienced by the capacitor.
Leakage current
For electrolytic capacitors, DC leakage current (DCL) is a special characteristic that other conventional capacitors do not have. This current is represented by the resistor Rleak in parallel with the capacitor in the series-equivalent circuit of electrolytic capacitors.
The reasons for leakage current are different between electrolytic capacitors with non-solid and with solid electrolyte or more common for "wet" aluminium and for "solid" tantalum electrolytic capacitors with manganese dioxide electrolyte as well as for electrolytic capacitors with polymer electrolytes. For non-solid aluminium electrolytic capacitors the leakage current includes all weakened imperfections of the dielectric caused by unwanted chemical processes taking place during the time without applied voltage (storage time) between operating cycles. These unwanted chemical processes depend on the kind of electrolyte. Water-based electrolytes are more aggressive to the aluminium oxide layer than are electrolytes based on organic liquids. This is why different electrolytic capacitor series specify different storage time without reforming.
Applying a positive voltage to a "wet" capacitor causes a reforming (self-healing) process which repairs all weakened dielectric layers, and the leakage current remain at a low level.
Although the leakage current of non-solid electrolytic capacitors is higher than current flow across the dielectric in ceramic or film capacitors, self-discharge of modern non-solid electrolytic capacitors with organic electrolytes takes several weeks.
The main causes of DCL for solid tantalum capacitors include electrical breakdown of the dielectric; conductive paths due to impurities or poor anodization; and bypassing of dielectric due to excess manganese dioxide, to moisture paths, or to cathode conductors (carbon, silver). This "normal" leakage current in solid electrolyte capacitors cannot be reduced by "healing", because under normal conditions solid electrolytes cannot provide oxygen for forming processes. This statement should not be confused with the self-healing process during field crystallization, see below, Reliability (Failure rate).
The specification of the leakage current in data sheets is often given as multiplication of the rated capacitance value CR with the value of the rated voltage UR together with an addendum figure, measured after a measuring time of two or five minutes, for example:
The leakage current value depends on the voltage applied, on the temperature of the capacitor, and on measuring time. Leakage current in solid MnO2 tantalum electrolytic capacitors generally drops very much faster than for non-solid electrolytic capacitors but remain at the level reached.
Dielectric absorption (soakage)
Dielectric absorption occurs when a capacitor that has remained charged for a long time discharges only incompletely when briefly discharged. Although an ideal capacitor would reach zero volts after discharge, real capacitors develop a small voltage from time-delayed dipole discharging, a phenomenon that is also called dielectric relaxation, "soakage" or "battery action".
Dielectric absorption may be a problem in circuits where very small currents are used in the function of an electronic circuit, such as long-time-constant integrators or sample-and-hold circuits. In most electrolytic capacitor applications supporting power supply lines, dielectric absorption is not a problem.
But especially for electrolytic capacitors with high rated voltage, the voltage at the terminals generated by the dielectric absorption can pose a safety risk to personnel or circuits. In order to prevent shocks, most very large capacitors are shipped with shorting wires that need to be removed before the capacitors are used.
Operational characteristics
Reliability (failure rate)
The reliability of a component is a property that indicates how reliably this component performs its function in a time interval. It is subject to a stochastic process and can be described qualitatively and quantitatively; it is not directly measurable. The reliability of electrolytic capacitors is empirically determined by identifying the failure rate in production accompanying endurance tests, see Reliability engineering.
Reliability normally is shown as a bathtub curve and is divided into three areas: early failures or infant mortality failures, constant random failures and wear out failures. Failures totalized in a failure rate are short circuit, open circuit, and degradation failures (exceeding electrical parameters).
The reliability prediction is generally expressed in a failure rate λ, abbreviated FIT (Failures In Time). This is the number of failures that can be expected in one billion (109) component-hours of operation (e.g., 1000 components for 1 million hours, or 1 million components for 1000 hours which is 1 ppm/1000 hours) at fixed working conditions during the period of constant random failures. This failure rate model implicitly assumes the idea of "random failure". Individual components fail at random times but at a predictable rate.
Billions of tested capacitor unit-hours would be needed to establish failure rates in the very low level range which are required today to ensure the production of large quantities of components without failures. This requires about a million units over a long time period, which means a large staff and considerable financing. The tested failure rates are often complemented with figures resulting from feedback from the field from major customers (field failure rate), which mostly results in a lower failure rate than tested.
The reciprocal value of FIT is Mean Time Between Failures (MTBF).
The standard operating conditions for FIT testing are 40 °C and 0.5 UR. For other conditions of applied voltage, current load, temperature, capacitance value, circuit resistance (for tantalum capacitors), mechanical influences and humidity, the FIT figure can be converted with acceleration factors standardized for industrial or military applications. The higher the temperature and applied voltage, the higher the failure rate, for example.
The most often cited source for failure rate conversion is MIL-HDBK-217F, the “bible” of failure rate calculations for electronic components. SQC Online, the online statistical calculator for acceptance sampling and quality control, provides an online tool for short examination to calculate given failure rate values for given application conditions.
Some manufacturers may have their own FIT calculation tables for tantalum capacitors. or for aluminium capacitors
For tantalum capacitors the failure rate is often specified at 85 °C and rated voltage UR as reference conditions and expressed as percent failed components per thousand hours (n %/1000 h). That is, “n” number of failed components per 105 hours, or in FIT the ten-thousand-fold value per 109 hours.
Tantalum capacitors are now very reliable components. Continuous improvement in tantalum powder and capacitor technologies have resulted in a significant reduction in the amount of impurities which formerly caused most field crystallization failures. Commercially available industrially produced tantalum capacitors now have reached as standard products the high MIL standard "C" level, which is 0.01%/1000 h at 85 °C and UR or 1 failure per 107 hours at 85 °C and UR. Converted to FIT with the acceleration factors coming from MIL HDKB 217F at 40 °C and 0.5 , UR is the failure rate. For a 100 μF/25 V tantalum chip capacitor used with a series resistance of 0.1 Ω the failure rate is 0.02 FIT.
Aluminium electrolytic capacitors do not use a specification in "% per 1000 h at 85 °C and UR". They use the FIT specification with 40 °C and 0.5 UR as reference conditions. aluminium electrolytic capacitors are very reliable components. Published figures show for low voltage types (6.3…160 V) FIT rates in the range of 1 to 20 FIT and for high voltage types (>160 …550 V) FIT rates in the range of 20 to 200 FIT. Field failure rates for aluminium e-caps are in the range of 0.5 to 20 FIT.
The published figures show that both tantalum and aluminium capacitor types are reliable components, comparable with other electronic components and achieving safe operation for decades under normal conditions. But a great difference exists in the case of wear-out failures. Electrolytic capacitors with non-solid electrolyte, have a limited period of constant random failures up to the point when wear-out failures begin. The constant random failure rate period corresponds to the lifetime or service life of “wet” aluminium electrolytic capacitors.
Lifetime
The lifetime, service life, load life or useful life of electrolytic capacitors is a special characteristic of non-solid aluminium electrolytic capacitors, whose liquid electrolyte can evaporate over time. Lowering the electrolyte level affects the electrical parameters of the capacitors. The capacitance decreases and the impedance and ESR increase with decreasing amounts of electrolyte. This very slow electrolyte drying-out depends on the temperature, the applied ripple current load, and the applied voltage. The lower these parameters compared to their maximum values, the longer the capacitor's “life”. The “end of life” point is defined by the appearance of wear-out failures or degradation failures when either capacitance, impedance, ESR or leakage current exceed their specified change limits.
The lifetime is a specification of a collection of tested capacitors and delivers an expectation of the behavior of similar types. This lifetime definition corresponds to the time of the constant random failure rate in the bathtub curve.
But even after exceeding the specified limits and the capacitors having reached their “end of life”, the electronic circuit is not in immediate danger; only the functionality of the capacitors is reduced. With today's high levels of purity in the manufacture of electrolytic capacitors it is not to be expected that short circuits occur after the end-of-life-point with progressive evaporation combined with parameter degradation.
The lifetime of non-solid aluminium electrolytic capacitors is specified in terms of “hours per temperature", like "2,000h/105 °C". With this specification the lifetime at operational conditions can be estimated by special formulas or graphs specified in the data sheets of serious manufacturers. They use different ways for specification, some give special formulas, others specify their e-caps lifetime calculation with graphs that consider the influence of applied voltage. The basic principle for calculating the time under operational conditions is the so-called “10-degree-rule”.
This rule is also known as the Arrhenius rule. It characterizes the change of thermic reaction speed. For every 10 °C lower temperature the evaporation is reduced by half. That means for every 10 °C reduction in temperature, the lifetime of capacitors doubles. If a lifetime specification of an electrolytic capacitor is, for example, 2000 h/105 °C, the capacitor's lifetime at 45 °C can be ”calculated” as 128,000 hours—that is roughly 15 years—by using the 10-degrees-rule.
However, solid polymer electrolytic capacitors, and aluminium, tantalum, and niobium electrolytic capacitors also have a lifetime specification. The polymer electrolyte exhibits a small deterioration of conductivity caused by thermal degradation of the conductive polymer. The electrical conductivity decreases as a function of time, in agreement with a granular metal type structure, in which aging is due to the shrinking of the conductive polymer grains. The lifetime of polymer electrolytic capacitors is specified in terms similar to non-solid electrolytic capacitors but its lifetime calculation follows other rules, leading to much longer operational lifetimes.
Tantalum electrolytic capacitors with solid manganese dioxide electrolyte do not have wear-out failures, so they do not have a lifetime specification in the sense of non-solid aluminium electrolytic capacitors. Also, tantalum capacitors with non-solid electrolyte, the "wet tantalums", do not have a lifetime specification because they are hermetically sealed.
Failure modes, self-healing mechanism and application rules
The many different types of electrolytic capacitors exhibit different electrical long-term behavior, intrinsic failure modes, and self-healing mechanisms. Application rules for types with an intrinsic failure mode are specified to ensure capacitors with high reliability and long life.
Performance after storage
All electrolytic capacitors are "aged" during manufacturing by applying the rated voltage at high temperature for a sufficient time to repair all cracks and weaknesses that may have occurred during production. However, a particular problem with non-solid aluminium models may occur after storage or unpowered periods. Chemical processes (corrosion) can weaken the oxide layer, which may lead to a higher leakage current. Most modern electrolytic systems are chemically inert and do not exhibit corrosion problems, even after storage times of two years or longer. Non-solid electrolytic capacitors using organic solvents like GBL as electrolyte do not have problems with high leakage current after prolonged storage. They can be stored for up to 10 years without problems
Storage times can be tested using accelerated shelf-life testing, which requires storage without applied voltage at the upper category temperature for a certain period, usually 1000 hours. This shelf life test is a good indicator for chemical stability and of the oxide layer, because all chemical reactions are accelerated by higher temperatures. Nearly all commercial series of non-solid electrolytic capacitors fulfill the 1000 hour shelf life test. However, many series are specified only for two years of storage. This also ensures solderability of the terminals.
For antique radio equipment or for electrolytic capacitors built in the 1970s or earlier, "preconditioning" may be appropriate. This is performed by applying the rated voltage to the capacitor via a series resistor of approximately 1 kΩ for one hour, allowing the oxide layer to repair itself through self-healing. Capacitors that fail leakage current requirements after preconditioning may have experienced mechanical damage.
Electrolytic capacitors with solid electrolytes do not have preconditioning requirements.
Causes of explosion
Electrolytic capacitors can explode due to several reasons, primarily related to internal pressure buildup and electrolyte issues:
Overvoltage and Reverse Polarity: Applying a voltage higher than the rated value or reversing the polarity can cause excessive current to flow, leading to rapid heating. This heating can decompose the electrolyte, generating gas and increasing internal pressure until the capacitor explodes.
Electrolyte Decomposition: Electrolytes in capacitors can decompose into gasses such as hydrogen when exposed to high temperatures or electrical stresses. The resulting gas increases internal pressure, leading to the rupture of the capacitor casing.
Design Flaws and Manufacturing Defects: Poor manufacturing processes can lead to weak points in the capacitor’s structure. During the capacitor plague, a widespread issue arose from the use of an incomplete electrolyte formula, which lacked necessary inhibitors to prevent gas formation and pressure buildup.
Thermal Runaway: High ripple currents can cause the capacitor to overheat. As the temperature rises, the electrolyte evaporates faster, creating more gas and increasing pressure, which can result in an explosion.
Aging and Deterioration: Over time, the materials within electrolytic capacitors degrade, reducing their ability to handle electrical stress and heat. This aging process can lead to internal failures and explosions.
Additional information
Capacitor symbols
Electrolytic capacitor symbols
Parallel connection
If an individual capacitor within a bank of parallel capacitors develops a short, the entire energy of the capacitor bank discharges through that short. Thus, large capacitors, particularly high voltage types, should be individually protected against sudden discharge.
Series connection
In applications where high withstanding voltages are needed, electrolytic capacitors can be connected in series. Because of individual variation in insulation resistance, and thus the leakage current when voltage is applied, the voltage is not distributed evenly across each series capacitor. This can result in the voltage rating of an individual capacitor being exceeded. A passive or active balancer circuit must be provided in order to equalize the voltage across each individual capacitor.
Polarity marking
Polarity marking for polymer electrolytic capacitors
Imprinted markings
Electrolytic capacitors, like most other electronic components, are marked, space permitting, with
manufacturer's name or trademark;
manufacturer's type designation;
polarity of the terminations (for polarized capacitors)
rated capacitance;
tolerance on rated capacitance
rated voltage and nature of supply (AC or DC)
climatic category or rated temperature;
year and month (or week) of manufacture;
certification marks of safety standards (for safety EMI/RFI suppression capacitors)
Smaller capacitors use a shorthand notation. The most commonly used format is: XYZ J/K/M “V”, where XYZ represents the capacitance (calculated as XY × 10Z pF), the letters K or M indicate the tolerance (±10% and ±20% respectively) and “V” represents the working voltage.
Examples:
105K 330V implies a capacitance of 10 × 105 pF = 1 μF (K = ±10%) with a rated voltage of 330 V.
476M 100V implies a capacitance of 47 × 106 pF = 47 μF (M = ±20%) with a rated voltage of 100 V.
Capacitance, tolerance and date of manufacture can be indicated with a short code specified in IEC/EN 60062. Examples of short-marking of the rated capacitance (microfarads): μ47 = 0,47 μF, 4μ7 = 4,7 μF, 47μ = 47 μF
The date of manufacture is often printed according to international standards.
Version 1: coding with year/week numeral code, "1208" is "2012, week number 8".
Version 2: coding with year code/month code. The year codes are: "R" = 2003, "S"= 2004, "T" = 2005, "U" = 2006, "V" = 2007, "W" = 2008, "X" = 2009, "A" = 2010, "B" = 2011, "C" = 2012, "D" = 2013, “E” = 2014 etc. Month codes are: "1" to "9" = Jan. to Sept., "O" = October, "N" = November, "D" = December. "X5" is then "2009, May"
For very small capacitors no marking is possible. Here only the traceability of the manufacturers can ensure the identification of a type.
Standardization
The standardization for all electrical, electronic components and related technologies follows the rules given by the International Electrotechnical Commission (IEC), a non-profit, non-governmental international standards organization.
The definition of the characteristics and the procedure of the test methods for capacitors for use in electronic equipment are set out in the Generic specification:
IEC/EN 60384-1 - Fixed capacitors for use in electronic equipment
The tests and requirements to be met by aluminium and tantalum electrolytic capacitors for use in electronic equipment for approval as standardized types are set out in the following sectional specifications:
IEC/EN 60384-3—Surface mount fixed tantalum electrolytic capacitors with manganese dioxide solid electrolyte
IEC/EN 60384-4—Aluminium electrolytic capacitors with solid (MnO2) and non-solid electrolyte
IEC/EN 60384-15—Fixed tantalum capacitors with non-solid and solid electrolyte
IEC/EN 60384-18—Fixed aluminium electrolytic surface mount capacitors with solid (MnO2) and non-solid electrolyte
IEC/EN 60384-24—Surface mount fixed tantalum electrolytic capacitors with conductive polymer solid electrolyte
IEC/EN 60384-25—Surface mount fixed aluminium electrolytic capacitors with conductive polymer solid electrolyte
IEC/EN 60384-26—Fixed aluminium electrolytic capacitors with conductive polymer solid electrolyte
Market
The market for electrolytic capacitors in 2008 was roughly 30% of the total market in value
aluminium electrolytic capacitors—US$3.9 billion (22%);
Tantalum electrolytic capacitors—US$2.2 billion (12%);
In number of pieces, these capacitors cover about 10% of the total capacitor market, or about 100 to 120 billion pieces.
Manufacturers and products
Date of the table: March 2015
| Technology | Components | null |
343857 | https://en.wikipedia.org/wiki/Lawn%20mower | Lawn mower | A lawn mower (also known as a grass cutter or simply mower, also often spelled lawnmower) is a device utilizing one or more revolving blades (or a reel) to cut a grass surface to an even height. The height of the cut grass may be fixed by the mower's design but generally is adjustable by the operator, typically by a single master lever or by a mechanism on each of the machine's wheels. The blades may be powered by manual force, with wheels mechanically connected to the cutting blades so that the blades spin when the mower is pushed forward, or the machine may have a battery-powered or plug-in electric motor. The most common self-contained power source for lawn mowers is a small 4-stroke (typically one-cylinder) internal combustion engine. Smaller mowers often lack any form of self-propulsion, requiring human power to move over a surface; "walk-behind" mowers are self-propelled, requiring a human only to walk behind and guide them. Larger lawn mowers are usually either self-propelled "walk-behind" types or, more often, are "ride-on" mowers that the operator can sit on and control. A robotic lawn mower ("lawn-mowing bot", "mowbot", etc.) is designed to operate either entirely on its own or less commonly by an operator on a remote control.
Two main styles of blades are used in lawn mowers. Lawn mowers employing a single blade that rotates about a single vertical axis are known as rotary mowers, while those employing a cutting bar and multiple blade assembly that rotates about a single horizontal axis are known as cylinder or reel mowers (although in some versions, the cutting bar is the only blade, and the rotating assembly consists of flat metal pieces which force the blades of grass against the sharp cutting bar).
There are several types of mowers, each suited to a particular scale and purpose. The smallest types, non-powered push mowers, are suitable for small residential lawns and gardens. Electrical or piston engine-powered push-mowers are used for larger residential lawns (although there is some overlap). Riding mowers, which sometimes resemble small tractors, are larger than push mowers and are suitable for large lawns. However, commercial riding lawn mowers (such as zero-turn mowers) can be "stand-on" types and often bear little resemblance to residential lawn tractors, being designed to mow large areas at high speed in the shortest time possible. The largest multi-gang (multi-blade) mowers are mounted on tractors and are designed for large expanses of grass such as golf courses and municipal parks, although they are ill-suited for complex terrain.
History
Invention
The lawn mower was invented in 1830 by Edwin Beard Budding of Stroud, Gloucestershire, England. Budding's mower was designed primarily to cut the grass on sports grounds and extensive gardens, as a superior alternative to the scythe, and was granted a British patent on August 31, 1830.
Budding's first machine was wide with a frame made of wrought iron. The mower was pushed from behind. Cast-iron gear wheels transmitted power from the rear roller to the cutting cylinder, allowing the rear roller to drive the knives on the cutting cylinder; the ratio was 16:1. Another roller placed between the cutting cylinder and the main or land roller could be raised or lowered to alter the height of cut. The grass clippings were hurled forward into a tray-like box. It was soon realized, however, that an extra handle was needed in front to help pull the machine along. Overall, these machines were remarkably similar to modern mowers.
Two of the earliest Budding machines sold went to Regent's Park Zoological Gardens in London and the Oxford colleges. In an agreement between John Ferrabee and Edwin Budding dated May 18, 1830, Ferrabee paid the costs of enlarging the small blades, obtained letters of patent and acquired rights to manufacture, sell and license other manufacturers in the production of lawn mowers. Without patent, Budding and Ferrabee were shrewd enough to allow other companies to build copies of their mower under licence, the most successful of these being Ransomes of Ipswich, which began making mowers as early as 1832.
His machine was the catalyst for the preparation of modern-style sporting ovals, playing fields (pitches), grass courts, etc. This led to the codification of modern rules for many sports, including for football, lawn bowls, lawn tennis and others.
Further improvements
It took ten more years and further innovations to create a machine that could be drawn by animals, and sixty years before a steam-powered lawn mower was built. In the 1850s, Thomas Green & Son of Leeds introduced a mower called the Silens Messor (meaning silent cutter), which used a chain drive to transmit power from the rear roller to the cutting cylinder. These machines were lighter and quieter than the gear-driven machines that preceded them, although they were slightly more expensive. The rise in popularity of lawn sports helped prompt the spread of the invention. Lawn mowers became a more efficient alternative to the scythe and domesticated grazing animals.
Manufacture of lawn mowers took off in the 1860s. By 1862, Ferrabee's company was making eight models in various roller sizes. He manufactured over 5000 machines until production ceased in 1863. The first grass boxes were flat trays but took their present shape in the 1860s. James Sumner of Lancashire patented the first steam-powered lawn mower in 1893. His machine burned petrol and/or paraffin (kerosene) as fuel. These were heavy machines that took several hours to warm up to operating pressure. After numerous advances, these machines were sold by the Stott Fertilizer and Insecticide Company of Manchester and Sumner. The company they both controlled was called the Leyland Steam Motor Company.
Around 1900, one of the best known English machines was the Ransomes' Automaton, available in chain- or gear-driven models. Numerous manufacturers entered the field with petrol (gasoline) engine-powered mowers after the start of the 20th century. In 1902, The first was produced by Ransomes. JP Engineering of Leicester, founded after World War I, produced a range of very popular chain-driven mowers. About this time, an operator could ride behind animals that pulled the large machines. These were the first riding mowers.
The first United States patent for a reel lawn mower was granted to Amariah Hills on January 12, 1868. In 1870, Elwood McGuire of Richmond, Indiana designed a human-pushed lawn mower, which was very lightweight and a commercial success. John Burr patented an improved rotary-blade lawn mower in 1899, with the wheel placement altered for better performance. Amariah Hills went on to found the Archimedean Lawn Mower Co. in 1871.
In the United States, gasoline-powered lawn mowers were first manufactured in 1914 by Ideal Power Mower Co. of Lansing, Michigan, based on a patent by Ransom E. Olds. Ideal Power Mower also introduced the world's first self-propelled, riding lawn tractor in 1922, known as the "Triplex". The roller-drive lawn mower has changed very little since around 1930. Gang mowers, those with multiple sets of blades to cut a wider swath, were built in the United States in 1919 by the Worthington Mower Company.
Atco Ltd and the first motor mower
In the 1920s one of the most successful companies to emerge during this period was Atco, at that time a brand name of Charles H Pugh Ltd. The Atco 'Standard' motor mower, launched in 1921 was an immediate success. Just 900 of the 22-inch-cut machines were made in 1921, each costing £75. Within five years, annual production had accelerated to tens of thousands. Prices were reduced and a range of sizes were available, making the Standard the first truly mass-produced engine-powered mower.
Rotary mowers
Rotary mowers were not developed until engines were small enough and powerful enough to run the blades at sufficient speed. Many people experimented with rotary blade mowers in the late 1920s and early 1930s, and Power Specialties Ltd. introduced a gasoline-powered rotary mower. Kut Kwick replaced the saw blade of the "Pulp Saw" with a double-edged blade and a cutter deck, converting the "Pulp Saw" into the first ever out-front rotary mower.
One company that produced rotary mowers commercially was the Australian Victa company, starting in 1952. Its mowers were lighter and easier to use than similar ones that had come before. The first Victa mowers were made at Mortlake, an inner suburb of Sydney, by local resident Mervyn Victor Richardson. He made his first model out of scrap in his garage. The first Victa mowers were then manufactured, going on sale on 20 September 1952. The new company, Victa Mowers Pty Ltd, was incorporated on 13 February 1953.
The venture was so successful that by 1958 the company moved to much larger premises in Parramatta Road, Concord, and then to Milperra, by which time the mower incorporated an engine, designed and manufactured by Victa, which was specially designed for mowing, rather than employing a general-purpose engine bought from outside suppliers. Two Victa mowers, from 1958 and 1968 respectively, are held in the collection of the National Museum of Australia. The Victa mower is regarded as something of an Australian icon, appearing en masse, in simulated form, at the opening of the Sydney Olympic Games in 2000.
The hover mower, first introduced by Flymo in 1964, is a form of rotary mower using an air cushion on the hovercraft principle.
Types
By rotation
Cylinder or reel mowers
A cylinder mower or reel mower carries a fixed, horizontal cutting blade at the desired height of cut. Over this is a fast-spinning reel of blades which force the grass past the cutting bar. Each blade in the blade cylinder forms a helix around the reel axis, and the set of spinning blades describes a cylinder.
Of all the mowers, a properly adjusted cylinder mower makes the cleanest cut of the grass, and this allows the grass to heal more quickly. The cut of a well-adjusted cylinder mower is straight and definite, as if cut with a pair of scissors. This clean cut promotes healthier, thicker and more resilient lawn growth that is more resistant to disease, weeds and parasites. Lawn cut with a cylinder mower is less likely to result in yellow, white or brown discolouration as a result of leaf shredding. While the cutting action is often likened to that of scissors, it is neither necessary nor desirable for the blades of the spinning cylinder to contact the horizontal cutting bar. When the reel touches the cutting bar the work required by the mower increases dramatically. When the gap between the blades is less than the thickness of the grass blades, a clean cut can still be made without additional friction. When the gap is greater than the thickness of the grass blades, grass will slip through the gap uncut. Reel mowers also have more difficulty mowing over uneven terrain.
There are many variants of the cylinder mower. Push mowers have no engine and are usually used on smaller lawn areas where access is a problem, where noise pollution is undesirable and where air pollution is unwanted. As the mower is pushed along, the wheels drive gears which rapidly spin the reel. Typical cutting widths are . Advances in materials and engineering have resulted in these mowers being very light and easy to operate and manoeuvre compared with their predecessors while still giving all the cutting advantages of professional cylinder mowers. Their distinct environmental benefits, both in noise and air pollution, are also strong selling points, something not lost on many international zoos, animal sanctuaries and exclusive hotel groups.
The basic push mower mechanism is also used in gangs towed behind a tractor. The individual mowers are arranged in a "v" behind the tractor with each mower's track slightly overlapping that of the mower in front of it. Gang mowers are used over large areas of turf such as sports fields or parks.
A gasoline engine or electric motor can be added to a cylinder mower to power the cylinder, the wheels, the roller, or any combination of these. A typical arrangement on electric powered machines for residential lawns is for the motor to power the cylinder while the operator pushes the mower along. The electric models can be corded or cordless. On petrol machines the engine drives both the cylinder and the rear roller. Some variants have only three blades in a reel spinning at great speed, and these models are able to cut grass which has grown too long for ordinary push mowers. One type of reel mower, now largely obsolete, was a powered version of the traditional side-wheel push mower, which was used on residential lawns. An internal combustion engine sat atop the reel housing and drove the wheels, usually through a belt. The wheels in turn drove the reel, as in the push mower.
Greens mowers are used for the precision cutting of golf greens and have a cylinder made up of at least eight, but normally ten, blades. The machine has a roller before and after the cutting cylinder which smooths the freshly cut lawn and minimizes wheel marks. Due to the weight, the engine also propels the mower. Much smaller and lighter variants of the roller mower are sometimes used for small patches of ornamental lawns around flower beds, and these have no engine.
Riding reel mowers are also produced. Typically, the cutting reels are ahead of the vehicle's main wheels, so that the grass can be cut before the wheels push the grass over onto the ground. The reels are often hydraulically powered.
The main parts of a cylinder or reel mower are:
Blade reel/cylinder: Consists of numerous (3 to 7) helical blades that are attached to a rotating shaft. The blades rotate, creating a scissor-like cutting motion against the bed knife.
Bed knife: The stationary cutting mechanism of a cylinder/reel mower. This is a fixed horizontal blade that is mounted to the frame of the mower.
Body frame: The main structural frame of the mower onto which the other parts of the mower are mounted.
Wheels: Help propel the mower in action. Generally, reel mowers have two wheels.
Push handle: The "power source" of a manually operated reel mower. This is a sturdy T-shaped, rectangular, or trapezoidal handle that is connected to the frame, wheels and blade chamber.
Motor: The power source of a reel mower that is powered by gasoline or electricity.
Rotary mowers
A rotary mower rotates about a vertical axis with the blade spinning at high speed relying on impact to cut the grass. This tends to result in a rougher cut and bruises and shreds the grass leaf resulting in discolouration of the leaf ends as the shredded portion dies. This is particularly prevalent if the blades become clogged or blunt. Most rotary mowers need to be set a little higher than cylinder equivalents to avoid scalping and gouging of slightly uneven lawns, although some modern rotaries are fitted with a rear roller to provide a more formal striped cut. These machines will also tend to cut lower () than a standard four-wheeled rotary.
The main parts of a rotary mower are:
Cutter deck housing: Houses the blade and the drive system of the mower. It is shaped to effectively eject the grass clippings from the mower.
Blade mounting and drive system: The blade of a rotary mower is usually mounted directly to the crankshaft of its engine, but it can be propelled by a hydraulic motor or a belt pulley system.
Mower blade: A blade that rotates in a horizontal plane (about a vertical axis). Some mowers have multiple blades. The blade features edges that are slightly curved upward to generate a continuous air flow as the blade rotates (as a fan), thus creating a sucking and tearing action.
Engine/motor: May be powered by gasoline or electricity.
Wheels: Generally four wheels, two front and two rear. Some mowers have a roller in place of the rear wheels.
By energy source
Gasoline (petrol)
Extensive grass trimming was not common before the widespread application of the vertical shaft single cylinder gasoline/petrol engine. In the United States this development paralleled the market penetration of companies such as the Briggs & Stratton company of Wisconsin.
Most rotary push mowers are powered by internal combustion engines. Such engines are usually four-stroke engines, used for their greater torque and cleaner combustion (although a number of older models used two-stroke engines), running on gasoline (petrol) or other liquid fuels. Internal combustion engines used with lawn mowers normally have only one cylinder. Power generally ranges from four to seven horsepower. The engines usually have a carburetor and require a manual pull crank to start them, although an electric starter is offered on some models, particularly large riding and commercial mowers. Some mowers have a throttle control on the handlebar with which the operator can adjust the engine speed. Other mowers have a fixed, pre-set engine speed. All are equipped with a governor (often centrifugal/mechanical or air vane style) to open the throttle as needed to maintain the pre-selected speed when the force needed to cut the thicker or taller grass is encountered. Gasoline mowers have the advantages over electric mowers of greater power and distance range. They do create a significant amount of pollution due to the combustion in the engine, and their engines require periodic maintenance such as cleaning or replacement of the spark plug and air filter, and changing the engine oil. California passed Assembly Bill 1356 an air pollution control law on October 9, 2021. The California bill barred sales of spark ignited (gasoline fueled) internal combustion engines less than 25HP used for farm or construction machines as of January 1, 2024. The California bill does not ban turf care machines larger than 25HP or those powered by compression ignition (diesel) engines.
Electricity
Electric mowers are further subdivided into corded and cordless electric models. Both are relatively quiet, typically producing less than 75 decibels, while a gasoline lawn mower can be 95 decibels or more.
Corded electric mowers are limited in range by their trailing power cord, which may limit their use with lawns extending outward more than 100–150 feet (30–45 m) from the nearest available power outlet. There is the additional hazard with these machines of accidentally mowing over the power cable, which stops the mower and may put users at risk of receiving a dangerous electric shock. Installing a residual-current device (GFCI) on the outlet may reduce the shock risk.
Cordless electric mowers are powered by a variable number (typically 1–4) of 12-to-80-volt rechargeable batteries. Typically, more batteries mean more run time and/or power (and more weight). Batteries can be in the interior of the lawnmower or on the outside. If on the outside, the depleted batteries can be quickly swapped with recharged batteries. Cordless mowers have the maneuverability of a gasoline-powered mower and the environmental friendliness of a corded electric mower, but they are more expensive and come in fewer models (particularly the self-propelling type) than either. The eventual disposal of worn-out batteries is problematic (though some manufacturers offer to recycle them), and the motors in some cordless mowers tend to be less powerful than gasoline motors of the same total weight (including batteries).
Propane
Lawn mowers powered by propane were also manufactured by Lehr.
By hand
In hand-powered lawn mowers, the reel is attached to the mower's wheels by gears, so that when the mower is pushed forward, the reel spins several times faster than the plastic or rubber-tired wheels turn. Depending on the placement of the reel, these mowers often cannot cut grass very close to lawn obstacles.
Other notable types
Hover mowers
Hover mowers are powered rotary push lawn mowers that use an impeller above the spinning blades to drive air downward, thereby creating an air cushion that lifts the mower above the ground. The operator can then easily move the mower as it floats over the grass. Hover mowers are necessarily light in order to achieve the air cushion and typically have plastic bodies with an electric motor. The most significant disadvantage, however, is the cumbersome usability in rough terrain or on the edges of lawns, as the lifting air-cushion is destroyed by wide gaps between the chassis and the ground. Hover mowers are built to operate on steep slopes, waterfronts, and high-weeded areas, so they are often used by golf course greenskeepers and commercial landscapers. Grass collection is often available, but can be poor in some models. The quality of cut can be inferior if the grass is pushed away from the blade by the cushion of air.
Robotic mowers
Tractor pulled mowers
Tractor pulled mowers are usually in the form of an attachment to a tractor. The attachments can simply function by the movement of the tractor similar to manual push cylinder mowers, but also sometimes may have powered moving blades. They are commonly mounted on either the side or the back of the tractor.
Riding lawn mowers
Riding mowers (U.S. and Canada) or ride-on mowers (U.K. and Canada) are a popular alternative for large lawns. The operator is provided with a seat and controls on the mower and literally rides on the machine. Most use the horizontal rotating blade system, though usually with multiple blades. A common form of ride-on mower is the lawn tractor. These are usually designed to resemble a small agricultural tractor, with the cutting deck mounted amidships between the front and rear axles.
The drives for these mowers are in several categories. The most common transmission for tractors is a manual transmission. The second most common transmission type is a form of continuously variable transmission, called hydrostatic transmission. These transmissions take several forms, from pumps driving separate motors, which may incorporate a gear reduction, to fully integrated units containing a pump, motor and gear reduction. Hydrostatic transmissions are more expensive than mechanical transmissions, but they are easier to use and can transmit greater torque to the wheels compared to a typical mechanical transmission. The least common drive type, and the most expensive, is electric.
There have been a number of attempts to replace hydrostatic transmissions with lower cost alternatives, but these attempts, which include variable belt types, e.g. MTD's "Auto Drive", and toroidal, have various performance or perception problems that have caused their market life to be short or their market penetration to be limited.
Riding lawn mowers can often mount other devices, such as rototillers/rotavators, snow plows, snow blowers, yard vacuums, occasionally even front buckets or fork-lift tines (these are more properly known as "lawn tractors" in this case, being designed for a number of tasks). The ability to tow other devices is because they have multiple gears, often up to 5 or 6 and variable top speeds. Compact utility tractors equipped with a belly mower can look similar to riding lawn mowers, but they are typically larger, equipped with diesel engines, and feature a three-point hitch and rollover protection structure; these features are generally absent on riding lawn mowers.
The deck of a rotary mower is typically made of steel. Lighter steel is used on less expensive models, and heavier steel on more expensive models for durability. Other deck materials include aluminium, which does not rust and is a staple of higher priced mowers, and hard composite plastic, which does not rust and is lighter and less expensive than aluminium. Electric mowers typically have a plastic deck.
Riding mowers typically have an opening in the side or rear of the housing where the cut grass is expelled, as do most rotary lawn mowers. Some have a grass catcher attachment at the opening to bag the grass clippings.
Mulching mowers
Mulching mowers use special mulching blades which are available for rotary mowers. The blade is designed to keep the clippings circulating underneath the mower until the clippings are chopped quite small. Other designs have twin blades to mulch the clippings to small pieces. This function has the advantages of forgoing the additional work collecting and disposing of grass clippings while reducing lawn waste in such a way that also creates convenient compost for the lawn, forgoing the expense and adverse environmental effect of fertilizer.
Mower manufacturers market their mowers as side discharge, 2-in-1, meaning bagging and mulching or side discharging and mulching, and 3-in-1, meaning bagging, mulching, and side discharge. Most 2-in-1 bagging and mulching mowers require a separate attachment to discharge grass onto the lawn. Some side discharge mower manufacturers also sell separate "mulching plates" that will cover the opening on the side discharge mower and, in combination with the proper blades, will convert the mower to a mulching mower. These conversions are impractical when compared with 2- or 3-in-1 mowers which can be converted in the field in seconds. There are two types of bagging mowers. A rear bag mower features an opening on the back of the mower through which the grass is expelled into the bag. Hi-vac mowers have a tunnel that extends from the side discharge to the bag. Hi-vac is also the type of grass collection used on some riding lawn mowers and lawn tractors and is suitable for use in dry conditions but less suitable for long wet lush grass as they often clog up. Mulching and bagging mowers are not well suited to long grass or thick weeds. In some ride-on mowers, the cut grass is dropped onto the ground and then collected by a set of rotating bristles, allowing even long, wet grass to be collected.
Rotary mowers with internal combustion engines come in three price ranges. Low priced mowers use older technology, smaller motors, and lighter steel decks. These mowers are targeted at the residential market and typically price is the most important selling point.
Professional mowers
Professional grass-cutting equipment, used by large establishments such as universities, sports stadiums and local authorities, usually take the form of much larger, dedicated, ride-on platforms or attachments that can be mounted on, or behind, a standard tractor unit (a "gang-mower"). Either type may use rotating-blade or cylindrical-blade type cutters, although high-quality mowed surfaces demand the latter. Wide-area mowers (WAMs) are commercial grade mowers which have decks extended to either side, many to . These extensions can be lowered for large area mowing or raised to decrease the mower's width and allow for easy transport on city roads or trailers. Commercial lawn-mowing companies have also enthusiastically adopted types such as the zero-turn mower (in both ride-on and stand-on versions), which allow high speed over the grass surface, and rapid turnaround at the end of rows, as well as excellent maneuverability around obstacles.
Mowers mounted on a tractor's three-point hitch may be known as finish mowers used for maintaining lawn, flail mowers used for maintaining rough grass on rough surfaces, or brush mowers used for cutting brush and small trees.
Safety issues
Rotary mowers can throw out debris with extreme velocity and energy. Additionally, the blades of a self-powered push mower (gasoline or electric) can injure a careless or inattentive user; consequently, many come equipped with a dead man's switch to immediately disable the blade rotation when the user is no longer holding the handle. In the United States, over 12,000 people per year are hospitalized as a result of lawn mower accidents. In 2016, 86,000 adults and 4,500 children were admitted to the emergency room for lawnmower injuries. The vast majority of these injuries can be prevented by wearing protective footwear when mowing. The American Academy of Pediatrics recommends that children be at least 12 years old before they are allowed to use a walk-behind lawn mower and at least 16 years of age before using a riding mower and that they "should not operate lawn mowers until they have displayed the necessary levels of judgment, strength, coordination, and maturity". Persons using a mower should wear heavy footwear, eye protection, and hearing protection in the case of engine-powered mowers.
Environmental and occupational impact
A 2001 study showed that some mowers produce the same amount of pollution (emissions other than carbon dioxide) in one hour as driving a 1992 model vehicle for . Another estimate puts the amount of pollution from a lawn mower at four times the amount from a car, per hour, although this report is no longer available. Beginning in 2011, the United States Environmental Protection Agency set standards for lawn equipment emissions and expects a reduction of at least 35 percent.
Gas powered lawn mowers produce GHG emissions. A minimum-maintained lawn management practice with clipping recycling, and minimum irrigation and mowing, is recommended to mitigate global warming effects from urban turfgrass system.
Battery-powered lawn mowers offer cleaner alternatives to consumers by producing zero emissions, being more efficient, and eliminating risks of spilled gasoline. Gasoline-powered lawnmowers are not regulated to have emission-capturing technology.
Mowers can create significant noise pollution, and could cause hearing loss if used without hearing protection for prolonged periods of time. Lawn mowers also present an occupational hearing hazard to the nearly one million people who work in lawn service and ground-keeping. One study assessed the occupational noise exposure among groundskeepers at several North Carolina public universities and found noise levels from push lawn mowers measured between 86 and 95 decibels (A-weighted) and from riding lawn mowers between 88 and 96 dB(A); both types exceeded the National Institute for Occupational Safety and Health (NIOSH) Recommended Exposure Limit of 85 dB(A).
The risk of hearing loss and noise pollution can be reduced by using battery-operated mowers or appropriate hearing protection such as earplugs or earmuffs.
It is possible for a lawn mower to damage the underlying soil, the roots of the grass, and the mower itself if the blades cut through the grass and collide with the underlying ground. Therefore, it is important to adjust mower height properly and choose the right tires for a lawn mower to prevent any marking or digging.
| Technology | Farm and garden machinery | null |
343884 | https://en.wikipedia.org/wiki/Lacerta%20%28genus%29 | Lacerta (genus) | Lacerta is a genus of lizards of the family Lacertidae.
Taxonomy
Lacerta was a fairly diverse genus containing around 40 species, until it was split into nine genera in 2007 by Arnold, Arribas & Carranza.
Fossil record
The earliest known members of the genus Lacerta are known from early Miocene epoch fossils indistinguishable in anatomy from the modern green lizards such as Lacerta viridis. Some fossil species from the ice-age mediterranean, such as Lacerta siculimelitensis, reached especially large sizes.
Species
The genus Lacerta contains the following species.
Some species formerly in Lacerta
Arranged alphabetically by specific name:
Anatololacerta anatolica – Anatolian rock lizard
Atlantolacerta andreanskyi – Atlas dwarf lizard, Andreansky's lizard
Iberolacerta aranica – Aran rock lizard
Iberolacerta aurelioi – Aurelio's rock lizard
Archaeolacerta bedriagae – Bedriaga's rock lizard
Iberolacerta bonnali – Pyrenean rock lizard
Apathya cappadocica – Anatolian lizard
Darevskia chlorogaster – Green-bellied lizard
Phoenicolacerta cyanisparsa
Omanosaura cyanura – blue-tailed Oman lizard
Anatololacerta danfordi – Danford's lizard
Darevskia defilippii – Elburs lizard
Darevskia dryada
Teira dugesii – Madeiran wall lizard
Geosaurus giganteus
Hellenolacerta graeca – Greek rock lizard
Iberolacerta horvathi – Horvath's rock lizard
Omanosaura jayakari – Jayakar's lizard
Phoenicolacerta kulzeri
Phoenicolacerta laevis – Lebanon lizard
Timon lepidus – ocellated lizard, foot lizard
Iberolacerta monticola – Iberian rock lizard
Dinarolacerta mosorensis - Mosor rock lizard
Podarcis muralis – common wall lizard
Anatololacerta oertzeni
Dalmatolacerta oxycephala – sharp-snouted rock lizard
Parvilacerta parva – dwarf lizard
Darevskia steineri
Zootoca vivipara – viviparous lizard
Iranolacerta zagrosica
| Biology and health sciences | Lizards and other Squamata | Animals |
343960 | https://en.wikipedia.org/wiki/Heavy%20bomber | Heavy bomber | Heavy bombers are bomber aircraft capable of delivering the largest payload of air-to-ground weaponry (usually bombs) and longest range (takeoff to landing) of their era. Archetypal heavy bombers have therefore usually been among the largest and most powerful military aircraft at any point in time. In the second half of the 20th century, heavy bombers were largely superseded by strategic bombers, which were often even larger in size, had much longer ranges and were capable of delivering nuclear bombs.
Because of advances in aircraft design and engineering — especially in powerplants and aerodynamics — the size of payloads carried by heavy bombers has increased at rates greater than increases in the size of their airframes. The largest bombers of World War I, the Zeppelin-Staaken Riesenflugzeuge of Germany, could carry a payload of up to of bombs; by the latter half of World War II, the Avro Lancaster (introduced in 1942) routinely delivered payloads of (and sometimes up to ) and had a range of , while the B-29 (1944) delivered payloads in excess of and had a range of . By the late 1950s, the jet-powered Boeing B-52 Stratofortress, travelling at speeds of up to (more than double that of a Lancaster), could deliver a payload of , over a combat radius of .
During World War II, mass production techniques made available large, long-range heavy bombers in such quantities as to allow strategic bombing campaigns to be developed and employed. This culminated in August 1945, when B-29s of the United States Army Air Forces dropped atomic bombs over Hiroshima and Nagasaki in Japan.
The arrival of nuclear weapons and guided missiles permanently changed the nature of military aviation and strategy. After the 1950s intercontinental ballistic missiles and ballistic missile submarines began to supersede heavy bombers in the strategic nuclear role. Along with the emergence of more accurate precision-guided munitions ("smart bombs") and nuclear-armed missiles, which could be carried and delivered by smaller aircraft, these technological advancements eclipsed the heavy bomber's once-central role in strategic warfare by the late 20th century. Heavy bombers have, nevertheless, been used to deliver conventional weapons in several regional conflicts since World War II (for example, B-52s in the Vietnam War).
Heavy bombers are now operated only by the air forces of the United States, Russia and China. They serve in both strategic and tactical bombing roles.
History
World War I
The first heavy bomber was designed as an airliner. Igor Sikorsky, an engineer educated in St Petersburg, but born in Kiev of Polish-Russian ancestry designed the Sikorsky Ilya Muromets to fly between his birthplace and his new home. It did so briefly until August 1914, when the Russo-Balt wagon factory converted to a bomber version, with British Sunbeam Crusader V8 engines in place of the German ones in the passenger plane. By December 1914 a squadron of 10 was bombing German positions on the Eastern Front and by summer 1916 there were twenty. It was well-armed with nine machine guns, including a tail gun and initially was immune to German and Austro-Hungarian air attack. The Sikorsky bomber had a wingspan just a few feet shorter than that of a World War II Avro Lancaster, while being able to carry a bomb load of only 3% of the later aircraft.
The Handley Page Type O/100 owed a lot to Sikorsky's ideas; of similar size, it used just two Rolls-Royce Eagle engines and could carry up to of bombs. The O/100 was designed at the beginning of the war for the Royal Navy specifically to sink the German High Seas Fleet in Kiel: the Navy called for “a bloody paralyser of an aircraft” Entering service in late 1916 and based near Dunkirk in France, it was used for daylight raids on naval targets, damaging a German destroyer. But after one was lost, the O/100 switched to night attacks.
The uprated Handley Page Type O/400 could carry a bomb, and wings of up to 40 were used by the newly formed, independent Royal Air Force from April 1918 to make strategic raids on German railway and industrial targets. A single O/400 was used to support T. E. Lawrence's Sinai and Palestine Campaign.
The Imperial German Air Service operated the Gotha bomber, which developed a series of marques. The Gotha G.IV operated from occupied Belgium from the Spring of 1917. It mounted several raids on London beginning in May 1917. Some reached no further than Folkestone or Sheerness on the Kent Coast. But on June 13, Gothas killed 162 civilians, including 18 children in a primary school, and injured 432 in East London. Initially, defence against air attack was poor, but by May 19, 1918, when 38 Gothas attacked London, six were shot down and another crashed on landing.
German aircraft companies also built a number of giant bombers, collectively known as the Riesenflugzeug. Most were produced in very small numbers from 1917 onwards and several never entered service. The most numerous were the Zeppelin-Staaken R.VI of which 13 saw service, bombing Russia and London: four were shot down and six lost on landing. The R.VIs were larger than the standard Luftwaffe bombers of World War II.
The Vickers Vimy, a long-range heavy bomber powered by two Rolls-Royce Eagle engines, was delivered to the newly formed Royal Air Force too late to see action (only one was in France at time of the Armistice with Germany). The Vimy's intended use was to bomb industrial and railway targets in western Germany, which it could reach with its range of and a bomb load of just over a ton. The Vickers Vimy is best known as the aircraft that made the first Atlantic crossing from St John's Newfoundland to Clifden in Ireland piloted by the Englishman John Alcock and navigated by Scot Arthur Whitten Brown on June 14, 1919.
Strategic bomber theory
Between the wars, aviation opinion fixed on two tenets. The first was that “the bomber will always get through.” The speed advantage of biplane fighters over bombers was insignificant, and it was believed that they would never catch them. Furthermore, there was no effective method of detecting incoming bombers at sufficiently long range to scramble fighters on an interception course. In practice, a combination of new radar technology and advances in monoplane fighter design eroded this disadvantage. Throughout the war, bombers continually managed to strike their targets, but suffered unacceptable losses in the absence of careful planning and escort fighters. Only the later de Havilland Mosquito light bomber was fast enough to evade fighters. Heavy bombers needed defensive armament for protection, which reduced their effective bomb payload.
The second tenet was that strategic bombing of industrial capacity, power generation, oil refineries, and coal mines could win a war. This was certainly vindicated by the firebombing of Japanese cities and the two atomic bombs dropped on Hiroshima and Nagasaki in August 1945, as Japan's fragile housing and cottage industry made themselves easily vulnerable to attack, thus completely destroying Japanese industrial production (see Air Raids on Japan). It was less evident that it held true for the bombing of Germany. During the war, German industrial production actually increased, despite a sustained Allied bombing campaign.
As the German Luftwaffe's main task was to support the army, it never developed a successful heavy bomber. The prime proponent of strategic bombing, Luftwaffe Chief of Staff General Walther Wever, died in an air crash in 1936 on the very day that the specification for the Ural bomber (later won by the Heinkel He 177 which saw only limited use against the Soviet Union and the United Kingdom) was published. After Wever's death, Ernst Udet, development director at the Air Ministry steered the Luftwaffe towards dive bombers instead.
World War II
When Britain and France declared war on Germany in September 1939, the RAF had no heavy bomber yet in service; heavy bomber designs had started in 1936 and ordered in 1938.
The Handley Page Halifax and Avro Lancaster both originated as twin-engine "medium" bombers, but were rapidly redesigned for four Rolls-Royce Merlin engines and rushed into service once the technical problems of the larger Rolls-Royce Vulture emerged in the Avro Manchester. The Halifax joined squadrons in November 1940 and flew its first raid against Le Havre on the night of 11–12 March 1941. British heavy bomber designs often had three gun turrets with a total of 8 machine guns.
In January 1941, the Short Stirling reached operational status and first combat missions were flown in February. It was based on the successful Short Sunderland flying boat and shared its Bristol Hercules radial engines, wing, and cockpit with a new fuselage. It carried up to of bombs—almost twice the load of a Boeing B-17 Flying Fortress—but over just a radius. Due to its thick, short wing it was able to out-turn the main German night fighters, the Messerschmitt Bf 110 and the Junkers Ju 88. Heavy bombers still needed defensive armament for protection, even at night. The Stirling's low operational ceiling of just —also caused by the thick wing—meant that it was usually picked on by night fighters; within five months, 67 of the 84 aircraft in service had been lost. The bomb bay layout limited the size and types of bombs carried and it was relegated to secondary duties such as tug and paratrooper transport.
Due to the absence of British heavy bombers, 20 United States Army Air Corps Boeing B-17 Flying Fortresses were lent to the RAF, which during July 1941 commenced daylight attacks on warships and docks at Wilhelmshaven and Brest. These raids were complete failures. After eight aircraft were lost due to combat or breakdown and with many engine failures, the RAF stopped daylight bombing by September. It was clear that the B-17C model was not combat ready and that its five machine guns provided inadequate protection.
Combat feedback enabled Boeing engineers to improve the aircraft; when the first model B-17E began operating from English airfields in July 1942, it had many more defensive gun positions including a vitally important tail gunner. Eventually, U.S. heavy bomber designs, optimized for formation flying, had 10 or more machine guns and/or cannons in both powered turrets and manually operated flexible mounts to deliver protective arcs of fire. These guns were located in tail turrets, side gun ports either just behind the bombardier's clear nose glazing as "cheek" positions, or midway along the rear fuselage sides as "waist" positions. U.S. bombers carried .50 caliber machine gun, and dorsal (spine/top of aircraft) and ventral (belly/bottom of aircraft) guns with powered turrets. All of these machine guns could defend against attack when beyond the range of fighter escort; eventually, a total of 13 machine guns were fitted in the B-17G model. In order to assemble combat boxes of several aircraft, and later combat wings formed of a number of boxes, assembly ships were used to speed up formation.
Even this extra firepower, which increased empty weight by 20% and required more powerful versions of the Wright Cyclone engine, was insufficient to prevent serious losses in daylight. Escort fighters were needed but the RAF interceptors such as the Supermarine Spitfire had very limited endurance. An early raid on Rouen-Sotteville rail yards in Brittany on August 17, 1942, required four Spitfire squadrons outbound and five more for the return trip.
The USAAF chose to attack aircraft factories and component plants. On August 17, 1943, 230 Fortresses attacked a ball-bearing plant in Schweinfurt and again two months later, with 291 bombers, in the second raid on Schweinfurt. The works was severely damaged but at a huge cost: 36 aircraft lost in the first raid, 77 in the second. Altogether 850 airmen were killed or captured; only 33 Fortresses returned from the October raid undamaged
With the arrival of North American P-51 Mustangs and the fitting of drop tanks to increase the range of the Republic P-47 Thunderbolt for the Big Week offensive, between February 20–25, 1944, bombers were escorted all the way to the target and back. Losses were reduced to 247 out of 3,500 sorties, still devastating but accepted at the time.
The Consolidated B-24 Liberator and later version of the Fortress carried even more extensive defensive armament fitted into Sperry ball turrets. This was a superb defensive weapon that rotated a full 360 degrees horizontally with a 90-degree elevation. Its twin M2 Browning machine guns had an effective range of . The Liberator was the result of a proposal to assemble Fortresses in Consolidated plants, with the company returning with its own design of a longer-range, faster and higher-flying aircraft that could carry an extra ton of bombs. Early orders were for France (delivered to the RAF after the fall of France) and Britain, already at war, with just a batch of 36 for the USAAF.
Neither the USAAF nor the RAF judged the initial design suitable for bombing and it was first used on a variety of VIP transport and maritime patrol missions. Its long range, however, persuaded the USAAF to send 177 Liberators from Benghazi in Libya to bomb the Romanian oilfields on August 1, 1943, in Operation Tidal Wave. Due to navigational errors and alerted German flak batteries and fighters, only half returned to base although a few landed safely at RAF bases in Cyprus and some in Turkey, where they were interned. Only 33 were undamaged. Damage to the refineries was soon repaired and oil production actually increased.
By October 1942, a new Ford Motor Company plant at Willow Run Michigan was assembling Liberators. Production reached a rate of over one an hour in 1944 helping the B-24 to become the most produced US aircraft of all time. It became the standard heavy bomber in the Pacific and the only one used by the RAAF. The SAAF used Liberators to drop weapons and ammunition during the Warsaw Uprising in 1944.
The Avro Manchester was a twin-engine bomber powered by the ambitious 24-cylinder Rolls-Royce Vulture, but was rapidly redesigned for four Rolls-Royce Merlin engines due to technical problems with the Vulture which caused the aircraft to be unreliable, under-powered and hastened its withdrawal from service. Reaching squadrons early in 1942, the redesigned bomber with four Merlin engines and longer wings was renamed Avro Lancaster; it could deliver a load of bombs or up to with special modifications. The Lancaster's bomb bay was undivided, so that bombs of extraordinary size and weight such as the 10-ton Grand Slam could be carried.
Barnes Wallis, deputy chief aircraft designer at Vickers, spent much time thinking about weapons that might shorten the war. He conceived his “Spherical Bomb, Surface Torpedo” after watching his daughter flip pebbles over water. Two versions of the 'bouncing bomb' were developed: the smaller Highball was to be used against ships and attracted essential British Admiralty funding for his project. A flying torpedo, of which half was Torpex torpedo explosive, it was developed specifically to sink the Tirpitz which was moored in Trondheim fjord behind torpedo nets. Development delays in the 'bouncing bomb' meant that another Barnes Wallis invention, the 5-ton Tallboy was deployed instead; two Tallboys dropped by Avro Lancasters from altitude hit at near-supersonic speed and capsized the Tirpitz on November 12, 1944. Upkeep, the larger version of the bouncing bomb, was used to destroy the Mohne and Eder dams by Lancasters from the specially recruited and trained No. 617 Squadron RAF, often known as "the Dam Busters", under Wing Commander Guy Gibson.
In March and April 1945, as the war in Europe was ending, Lancasters dropped Grand Slams and Tallboys on U-boat pens and railway viaducts across north Germany. At Bielefeld more than of railway viaduct was destroyed by Grand Slams creating an earthquake effect, which shook the foundations.
The Boeing B-29 Superfortress was a development of the Fortress, but a larger design with four Wright R-3350 Duplex-Cyclone engines of much greater power, enabling it to fly higher, faster, further and with a bigger bomb load. The mammoth new Wright radial engines were susceptible to overheating if anything malfunctioned, and technical problems with the powerplant seriously delayed the B-29's operational service debut. The aircraft had four remotely operated twin-gun turrets on its fuselage, controlled through an analog computer sighting system; the operator could use any of a trio of Perspex ball stations. Only the tail gunner manually controlled his gun turret station in the rear of the airplane.
B-29s were initially deployed to bases in India and China, from which they could reach Japan; but the logistics (including transport of fuel for the B-29 fleet over the Himalayan range) of flying from these remote, primitive airfields were complicated and costly. The island of Saipan in the Marianas was assaulted to provide Pacific air bases from which to bomb Japanese cities. Initial high-level, daylight bombing raids using high-explosive bombs on Japanese cities with their wood and paper houses produced disappointing results; the bombers were then switched to low-level, nighttime incendiary attacks for which they had not originally been designed (one variant, the B-29B was specially modified for low altitude night missions by removal of armament and other equipment). Japan burned furiously from the B-29 incendiary raids. On August 6, 1945, B-29 Enola Gay dropped an atomic bomb on Hiroshima. Three days later, B-29 Bockscar dropped another on Nagasaki. The war ended when Japan announced its surrender to the Allies on August 15, and the Japanese government subsequently signed the official instrument of surrender on September 2, 1945.
After World War II
After World War II, the name strategic bomber came into use, for aircraft that could carry aircraft ordnances over long distances behind enemy lines. They were supplemented by smaller fighter-bombers with less range and lighter bomb load, for tactical strikes. Later these were called strike fighters, attack aircraft and multirole combat aircraft.
When North Korea attacked South Korea in 1950 the USAF responded with daylight bomber raids on supply lines through North Korea. B-29 Superfortresses flew from Japan on behalf of the United Nations, but the supply line for North Korea's army from the Soviet Union was physically and politically out of reach: North Korea for the most part lacked worthwhile strategic targets of its own. The Soviet-backed Northern forces easily routed the South Korean army. The distance to North Korea was too great for fighter escorts based in Japan, so the B-29s flew alone. In November, Mikoyan-Gurevich MiG-15s flown by Soviet pilots started to intercept the US bombers over North Korea. The MiG-15 was specifically designed to destroy US heavy bombers; it could out-perform any fighter deployed by United Nations air forces until the capable F-86 Sabre was produced in greater numbers and brought to Korea. After 28 B-29s were lost, the bombers were restricted to night interdiction and concentrated on destroying supply routes, including the bridges over the Yalu river into China.
By the 1960s, manned heavy bombers could not match the intercontinental ballistic missile in the strategic nuclear role. More accurate precision-guided munitions ("smart bombs"), nuclear-armed missiles or bombs were able to be carried by smaller aircraft such as fighter-bombers and multirole fighters. Despite these technological innovations and new capabilities of other contemporary military aircraft, large strategic bombers such as the B-1, B-52 and B-2 have been retained for the role of carpet bombing in several conflicts. The most prolific example (in terms of total bomb tonnage) is the U.S. Air Force B-52 Stratofortress during the 1960s–early 1970s Vietnam War era, in Operation Menu, Operation Freedom Deal, and Operation Linebacker II. In 1987 the Soviet Tu-160—the heaviest supersonic bomber/aircraft currently in active service—entered service; it can carry twelve long-range cruise missiles.
The 2010 New START agreement between the United States of America and the Russian Federation defined a "heavy bomber" by two characteristics:
range greater than
equipped for long-range nuclear "air-launched cruise missiles" (ALCMs), defined as an air-to-surface cruise missile of a type flight-tested from an aircraft or deployed on a bomber after 1986.
List of heavy bombers
Some notable heavy bombers are listed below
World War I
AEG G.I, II, III, IV, V, R.I
Airco DH.10 Amiens
Albatros G.I, II, III
Blériot 73
Caproni Ca.1, Ca.3 Ca.4, Ca.5
Caudron R.11
Farman F.50
Friedrichshafen G.II, G.III, G.IV
Gotha G.I, G.II, G.III, G.IV G.V
Handley Page Type O/100 & O/400
Handley Page V/1500
Letord Let.3 6 & 7
Martin MB-1
Rumpler G.I, G.II & G.III
Short Bomber
Sikorsky Ilya Muromets
Vickers Vimy
Zeppelin-Staaken R.VI
Zeppelin-Staaken R.XV
Zeppelin-Staaken R.XVI
Interwar period (interbellum)
Armstrong Whitworth Whitley
Avro 549 Aldershot
Avro Manchester
Blériot 127
Boeing B-17 Flying Fortress
Bréguet 410
Bristol Bombay
Curtiss B-2 Condor
Dornier Do 11
Douglas Y1B-7
Douglas B-18 Bolo
Farman BN.4
Farman F.60 Goliath
Farman F.120
Farman F.220
Fokker XB-8
Friedrichshafen G.V
Handley Page H.P.54 Harrow
Handley Page Heyford
Handley Page Hinaidi
Handley Page Hyderabad
Lioré et Olivier LeO 20
Martin NBS-1/MB-2
Mitsubishi Ki-20
Mitsubishi Ki-21
Petlyakov Pe-8
Tupolev TB-1
Tupolev TB-3
Vickers Virginia
World War II
Avro Lancaster
Avro Lincoln
Boeing B-17 Flying Fortress
Consolidated B-24 Liberator
Boeing B-29 Superfortress
Consolidated B-32 Dominator
Handley Page Halifax
Heinkel He 177
Piaggio P.108
Savoia-Marchetti SM.82
Short Stirling
Vickers Warwick
Post-WW2
Avro Vulcan
Boeing B-47 Stratojet
Boeing B-50 Superfortress
Boeing B-52 Stratofortress
Convair B-36 Peacemaker
Convair B-58 Hustler
Handley Page Victor
Ilyushin Il-28
Myasishchev M-4
Northrop Grumman B-2 Spirit
Northrop Grumman B-21 Raider
Rockwell B-1 Lancer
Short Sperrin
Sukhoi T-4
Tupolev Tu-4
Tupolev Tu-16
Tupolev Tu-22
Tupolev Tu-22M
Tupolev Tu-95
Tupolev Tu-160
Vickers Valiant
Xian H-6
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343971 | https://en.wikipedia.org/wiki/Medium%20bomber | Medium bomber | A medium bomber is a military bomber aircraft designed to operate with medium-sized bombloads over medium range distances; the name serves to distinguish this type from larger heavy bombers and smaller light bombers. Mediums generally carried about two tons of bombs, compared to light bombers that carried one ton, and heavies that carried four or more.
The term was used prior to and during World War II, based on available parameters of engine and aeronautical technology for bomber aircraft designs at that time. After the war, medium bombers were replaced in world air forces by more advanced and capable aircraft.
History
In the early 1930s many air forces were looking to modernize their existing bomber aircraft fleets, which frequently consisted of older biplanes. The new designs were typically twin-engined monoplanes, often of all-metal construction, and optimized for high enough performance and speed to help evade rapidly evolving fighter aircraft designs of the time. Some of these bombers, such as the Heinkel He 111, Junkers Ju 86, Savoia-Marchetti SM.79, Douglas B-18, and Armstrong Whitworth Whitley were developed from or in conjunction with existing airliners or transport aircraft.
The World War II-era medium bomber was generally considered to be any level bomber design that delivered about of ordnance over ranges of about . Typical heavy bombers were those with a nominal load of or more, and light bombers carried up to 2,000 lb (907 kg).
These distinctions were beginning to disappear by the middle of World War II, when the average fighter aircraft could now carry a 2,000 lb (907 kg) bombload. Advances in powerplants and designs eventually allowed light bombers, tactical bombers, and later jet fighter-bombers to take over the roles performed by mediums.
After the war, use of the term generally vanished; some of this was due to mass demobilization of the participant air forces' existing equipment, and the fact that several of the most-produced medium bomber types were now technologically obsolescent. Although a number of later aircraft were designed in this performance and load-carrying range, they were henceforth referred to as tactical bombers or strike aircraft instead. Examples of post-war mediums include the English Electric Canberra (along with its derived U.S. counterpart, the Martin B-57) and the Soviet Ilyushin Il-28 "Beagle".
Subsequent to World War II, only the U.S. Strategic Air Command ever used the term "medium bomber" in the 1950s to distinguish its Boeing B-47 Stratojets from somewhat larger contemporary Boeing B-52 Stratofortress "heavy bombers" in bombardment wings (older B-29 and B-50 heavy bombers were also redesignated as "medium" during this period). This nomenclature was purely semantic and bureaucratic, however as both the B-47 and B-52 strategic bombers were much larger and had far greater performance and load-carrying ability than any of the World War II-era heavy or medium bombers. Similarly, the Royal Air Force referred at times to its V bomber force as medium bombers, but this was in terms of range rather than load-carrying capacity.
Although the term is no longer used, development of aircraft that fulfil a 'medium bomber' mission in all but name continued and these have been employed in various post-World War II conflicts; examples include dedicated tactical bombers such as the Su-24, Su-34, F-111, J-16 and F-15E which have greater payload and range capability than fighter-bombers, but less than heavier strategic bombers.
Medium bombers
Introduced prior to World War II (September 1, 1939)
Armstrong Whitworth Whitley — first of three British medium bombers
Bloch MB.210
CANT Z.1007
Douglas B-18 Bolo — developed from the DC-2 airliner design
Douglas B-23 Dragon
Dornier Do 23
Fiat BR.20 Cicogna — first all-metal Italian bomber
Fokker T.V — Dutch army air force (Luchtvaartafdeling) bomber
Handley Page Hampden — British medium bomber, almost as fast as the Bristol Blenheim
Heinkel He 111 — considered a heavy bomber by the Luftwaffe for some missions
Ilyushin DB-3 — precursor to the Il-4 (see below)
Junkers Ju 52 (briefly during the Spanish Civil War)
Junkers Ju 86
Lioré et Olivier LeO 45 — fast French medium bomber
Martin B-10 — American bomber which was highly advanced at the time of its 1934 service introduction
Mitsubishi G3M — known to the Allies as "Nell"
Mitsubishi Ki-21 — "Sally"; replaced some Fiat BR.20 bombers in the Imperial Japanese Army Air Service
PZL.37 Łoś - the most advanced Polish aircraft at the time of the invasion of Poland
Savoia-Marchetti SM.79 — three-engined Italian medium bomber used successfully as a torpedo bomber early in World War II
Vickers Wellington — most-produced British medium bomber, with a unique aluminum lattice airframe designed by Barnes Wallis and capable of 2,500 miles range
World War II
de Havilland Mosquito — considered a multi-role aircraft
Dornier Do 217 — considered a heavy bomber by the Luftwaffe for some missions
Ilyushin Il-4 — long ranged Soviet bomber
Junkers Ju 88 — versatile aircraft used in many different roles, including torpedo bomber, dive bomber, night fighter and reconnaissance
Martin B-26 Marauder — had lowest mission loss rate of any USAAF bomber in World War II
Mitsubishi G4M — known to the Allies as "Betty"
Mitsubishi Ki-67 Hiryū — Allied reporting name "Peggy"; classified as heavy by the Imperial Japanese Army Air Service
Nakajima Ki-49 Donryu — "Helen"
North American B-25 Mitchell — most-produced American medium bomber
Savoia-Marchetti SM.84 — less successful replacement for the SM.79, which instead remained in service even after this aircraft
Tupolev Tu-2 — used in multiple roles similar to the German Junkers Ju 88
Yermolayev Yer-2 - developed from the failed Bartini Stal-7; used both as long-range bomber and as heavy ground attacker, with the latter role resulting in heavy losses.
Yokosuka P1Y Ginga — a medium bomber to the Imperial Japanese Navy Air Service; but in size, weight, speed etc. similar to Allied light bombers such as the Douglas A-26 Invader
Post World War II
English Electric Canberra - British jet bomber introduced in the 1950s
Ilyushin Il-28 — Soviet jet bomber
Martin B-57 Canberra — U.S. licence-built development of English Electric Canberra
Douglas A-3 Skywarrior — U.S. naval jet strike aircraft; a derivative was adopted by United States Air Force Tactical Air Command designated as the B-66 "tactical light bomber"
North American AJ Savage — U.S. naval medium bomber powered by two piston engines and a turbojet buried in the rear fuselage.
North American A-5 Vigilante — U.S. naval supersonic jet strike aircraft
Dassault Mirage IV — French supersonic jet strategic bomber
Grumman A-6 Intruder — U.S. naval strike aircraft; approximate in size to a World War II medium with range and payload comparable to a World War II heavy bomber
Blackburn Buccaneer — Royal Navy carrier-based maritime strike aircraft
General Dynamics F-111 — U.S. supersonic swing wing strike aircraft
Sukhoi Su-24 — Soviet supersonic swing wing strike aircraft similar in role and configuration to the F-111
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344195 | https://en.wikipedia.org/wiki/Living%20fossil | Living fossil | A living fossil is a deprecated term for an extant taxon that phenotypically resembles related species known only from the fossil record. To be considered a living fossil, the fossil species must be old relative to the time of origin of the extant clade. Living fossils commonly are of species-poor lineages, but they need not be. While the body plan of a living fossil remains superficially similar, it is never the same species as the remote relatives it resembles, because genetic drift would inevitably change its chromosomal structure.
Living fossils exhibit stasis (also called "bradytely") over geologically long time scales. Popular literature may wrongly claim that a "living fossil" has undergone no significant evolution since fossil times, with practically no molecular evolution or morphological changes. Scientific investigations have repeatedly discredited such claims.
The minimal superficial changes to living fossils are mistakenly declared as an absence of evolution, but they are examples of stabilizing selection, which is an evolutionary process—and perhaps the dominant process of morphological evolution.
The term is currently deprecated among paleontologists and evolutionary biologists.
Characteristics
Living fossils have two main characteristics, although some have a third:
Living organisms that are members of a taxon that has remained recognizable in the fossil record over an unusually long time span.
They show little morphological divergence, whether from early members of the lineage, or among extant species.
They tend to have little taxonomic diversity.
The first two are required for recognition as a living fossil; some authors also require the third, others merely note it as a frequent trait.
Such criteria are neither well-defined nor clearly quantifiable, but modern methods for analyzing evolutionary dynamics can document the distinctive tempo of stasis. Lineages that exhibit stasis over very short time scales are not considered living fossils; what is poorly-defined is the time scale over which the morphology must persist for that lineage to be recognized as a living fossil.
The term living fossil is much misunderstood in popular media in particular, in which it often is used meaninglessly. In professional literature the expression seldom appears and must be used with far more caution, although it has been used inconsistently.
One example of a concept that could be confused with "living fossil" is that of a "Lazarus taxon", but the two are not equivalent; a Lazarus taxon (whether a single species or a group of related species) is one that suddenly reappears, either in the fossil record or in nature, as if the fossil had "come to life again". In contrast to "Lazarus taxa", a living fossil in most senses is a species or lineage that has undergone exceptionally little change throughout a long fossil record, giving the impression that the extant taxon had remained identical through the entire fossil and modern period. Because of the mathematical inevitability of genetic drift, though, the DNA of the modern species is necessarily different from that of its distant, similar-looking ancestor. They almost certainly would not be able to cross-reproduce, and are not the same species.
The average species turnover time, meaning the time between when a species first is established and when it finally disappears, varies widely among phyla, but averages about 2–3million years. A living taxon that had long been thought to be extinct could be called a Lazarus taxon once it was discovered to be still extant. A dramatic example was the order Coelacanthiformes, of which the genus Latimeria was found to be extant in 1938. About that there is little debate – however, whether Latimeria resembles early members of its lineage sufficiently closely to be considered a living fossil as well as a Lazarus taxon has been denied by some authors in recent years.
Coelacanths disappeared from the fossil record some 80million years ago (in the upper Cretaceous period) and, to the extent that they exhibit low rates of morphological evolution, extant species qualify as living fossils. It must be emphasised that this criterion reflects fossil evidence, and is totally independent of whether the taxa had been subject to selection at all, which all living populations continuously are, whether they remain genetically unchanged or not.
This apparent stasis, in turn, gives rise to a great deal of confusion – for one thing, the fossil record seldom preserves much more than the general morphology of a specimen. To determine much about its physiology is seldom possible; not even the most dramatic examples of living fossils can be expected to be without changes, no matter how persistently constant their fossils and the extant specimens might seem. To determine much about noncoding DNA is hardly ever possible, but even if a species were hypothetically unchanged in its physiology, it is to be expected from the very nature of the reproductive processes, that its non-functional genomic changes would continue at more-or-less standard rates. Hence, a fossil lineage with apparently constant morphology need not imply equally constant physiology, and certainly neither implies any cessation of the basic evolutionary processes such as natural selection, nor reduction in the usual rate of change of the noncoding DNA.
Some living fossils are taxa that were known from palaeontological fossils before living representatives were discovered. The most famous examples of this are:
Coelacanthiform fishes (2 species)
Metasequoia, the dawn redwood discovered in a remote Chinese valley (1 species)
Glypheoid lobsters (2 species)
Mymarommatid wasps (10 species)
Eomeropid scorpionflies (1 species)
Jurodid beetles (1 species)
Soft sea urchins (59 species)
All the above include taxa that originally were described as fossils but now are known to include still-extant species.
Other examples of living fossils are single living species that have no close living relatives, but are survivors of large and widespread groups in the fossil record. For example:
Ginkgo biloba
Syntexis libocedrii, the cedar wood wasp
Dinoflagellates (typified on coccoid dinocysts: occasionally calcareous cell remnants)
All of these were described from fossils before later being found alive.
The fact that a living fossil is a surviving representative of an archaic lineage does not imply that it must retain all the "primitive" features (plesiomorphies) of its ancestral lineage. Although it is common to say that living fossils exhibit "morphological stasis", stasis, in the scientific literature, does not mean that any species is strictly identical to its ancestor, much less remote ancestors.
Some living fossils are relicts of formerly diverse and morphologically varied lineages, but not all survivors of ancient lineages necessarily are regarded as living fossils. See for example the uniquely and highly autapomorphic oxpeckers, which appear to be the only survivors of an ancient lineage related to starlings and mockingbirds.
Evolution and living fossils
The term living fossil is usually reserved for species or larger clades that are exceptional for their lack of morphological diversity and their exceptional conservatism, and several hypotheses could explain morphological stasis on a geologically long time-scale. Early analyses of evolutionary rates emphasized the persistence of a taxon rather than rates of evolutionary change. Contemporary studies instead analyze rates and modes of phenotypic evolution, but most have focused on clades that are thought to be adaptive radiations rather than on those thought to be living fossils. Thus, very little is presently known about the evolutionary mechanisms that produce living fossils or how common they might be. Some recent studies have documented exceptionally low rates of ecological and phenotypic evolution despite rapid speciation. This has been termed a "non-adaptive radiation" referring to diversification not accompanied by adaptation into various significantly different niches. Such radiations are explanation for groups that are morphologically conservative. Persistent adaptation within an adaptive zone is a common explanation for morphological stasis. The subject of very low evolutionary rates, however, has received much less attention in the recent literature than that of high rates.
Living fossils are not expected to exhibit exceptionally low rates of molecular evolution, and some studies have shown that they do not. For example, on tadpole shrimp (Triops), one article notes, "Our work shows that organisms with conservative body plans are constantly radiating, and presumably, adapting to novel conditions... I would favor retiring the term 'living fossil' altogether, as it is generally misleading." Some scientists instead prefer a new term stabilomorph, being defined as "an effect of a specific formula of adaptative strategy among organisms whose taxonomic status does not exceed genus-level. A high effectiveness of adaptation significantly reduces the need for differentiated phenotypic variants in response to environmental changes and provides for long-term evolutionary success."
The question posed by several recent studies pointed out that the morphological conservatism of coelacanths is not supported by paleontological data. In addition, it was shown recently that studies concluding that a slow rate of molecular evolution is linked to morphological conservatism in coelacanths are biased by the a priori hypothesis that these species are 'living fossils'. Accordingly, the genome stasis hypothesis is challenged by the recent finding that the genome of the two extant coelacanth species L. chalumnae and L. menadoensis contain multiple species-specific insertions, indicating transposable element recent activity and contribution to post-speciation genome divergence. Such studies, however, challenge only a genome stasis hypothesis, not the hypothesis of exceptionally low rates of phenotypic evolution.
History
The term was coined by Charles Darwin in his On the Origin of Species from 1859, when discussing Ornithorhynchus (the platypus) and Lepidosiren (the South American lungfish):
Other definitions
Long-enduring
A living taxon that lived through a large portion of geologic time.
The Australian lungfish (Neoceratodus fosteri), also known as the Queensland lungfish, is an example of an organism that meets this criterion. Fossils identical to modern specimens have been dated at over 100million years old. Modern Queensland lungfish have existed as a species for almost 30million years. The contemporary nurse shark has existed for more than 112million years, making this species one of the oldest, if not actually the oldest extant vertebrate species.
Resembles ancient species
A living taxon morphologically and/or physiologically resembling a fossil taxon through a large portion of geologic time (morphological stasis).
Retains many ancient traits
A living taxon with many characteristics believed to be primitive. This is a more neutral definition. However, it does not make it clear whether the taxon is truly old, or it simply has many plesiomorphies. Note that, as mentioned above, the converse may hold for true living fossil taxa; that is, they may possess a great many derived features (autapomorphies), and not be particularly "primitive" in appearance.
Relict population
Any one of the above three definitions, but also with a relict distribution in refuges.
Some paleontologists believe that living fossils with large distributions (such as Triops cancriformis) are not real living fossils. In the case of Triops cancriformis (living from the Triassic until now), the Triassic specimens lost most of their appendages (mostly only carapaces remain), and they have not been thoroughly examined since 1938.
Low diversity
Any of the first three definitions, but the clade also has a low taxonomic diversity (low diversity lineages).
Oxpeckers are morphologically somewhat similar to starlings due to shared plesiomorphies, but are uniquely adapted to feed on parasites and blood of large land mammals, which has always obscured their relationships. This lineage forms part of a radiation that includes Sturnidae and Mimidae, but appears to be the most ancient of these groups. Biogeography strongly suggests that oxpeckers originated in eastern Asia and only later arrived in Africa, where they now have a relict distribution.
The two living species thus seem to represent an entirely extinct and (as Passerida go) rather ancient lineage, as certainly as this can be said in the absence of actual fossils. The latter is probably due to the fact that the oxpecker lineage never occurred in areas where conditions were good for fossilization of small bird bones, but of course, fossils of ancestral oxpeckers may one day turn up enabling this theory to be tested.
Operational definition
An operational definition was proposed in 2017, where a 'living fossil' lineage has a slow rate of evolution and occurs close to the middle of morphological variation (the centroid of morphospace) among related taxa (i.e. a species is morphologically conservative among relatives). The scientific accuracy of the morphometric analyses used to classify tuatara as a living fossil under this definition have been criticised however, which prompted a rebuttal from the original authors.
Examples
Some of these are informally known as "living fossils".
Bacteria
Cyanobacteria – the oldest living fossils, emerging 3.5 billion years ago. They exist as single bacteria or in the form of stromatolites, layered rocks produced by colonies of cyanobacteria.
Protists
The dinoflagellate †Calciodinellum operosum.
The dinoflagellate †Dapsilidinium pastielsii.
The dinoflagellate †Posoniella tricarinelloides.
The coccolithophore Tergestiella adriatica.
Plants
Moss
Pteridophytes
Horsetails – Equisetum
Lycopods
Tree ferns and ferns
Gymnosperms
Conifers
Agathis – kauri in New Zealand, Australia and the Pacific and almasiga in the Philippines
Araucaria araucana – the monkey puzzle tree (as well as other extant Araucaria species)
Metasequoia – dawn redwood (Cupressaceae; related to Sequoia and Sequoiadendron)
Sciadopitys – a unique conifer endemic to Japan known in the fossil record for about 230 million years.
Taiwania cryptomerioides – one of the largest tree species in Asia.
Wollemia tree (Araucariaceae – a borderline example, related to Agathis and Araucaria)
Cycads – although this has been challenged by multiple lines of evidence
Ginkgo tree (Ginkgoaceae)
Welwitschia
Angiosperms
Amborella – a plant from New Caledonia, possibly closest to base of the flowering plants
Magnolia – a genus whose form is little changed since the earliest days of flowering plant evolution in the Cretaceous and possibly earlier
Trapa – water caltrops, seeds, and leaves of numerous extinct species are known all the way back to the Cretaceous.
Nelumbo – several species of lotus flower are known exclusively from fossils dating back to the Cretaceous.
Sassafras – many fossils of sassafras are known from the late cretaceous through the late pleistocene.
Platanus Sycamore fossils are very abundant throughout the northern hemisphere with several extinct species. Sycamore leaves and fruits are quite common in plant fossils. Sycamores exhibit many primitive features as well such their exfoliating bark which is a result of a lack of elasticity. Platanus Occidentalis fossils are known from the pliocene and the pleistocene in North America.
Nyssa Blackgum fossils go way back to the late cretaceous period. Many extinct species are recorded as well.
Liriodendron Fossils from the cretaceous and the tertiary period are found with many extinct species. Tulip Trees at one point were present in europe during the cretaceous and the early paleocene. Liriodendron Tulipifera fossils dating from the pliocene and pleistocene were discovered at the chowan formation in North Carolina.
Liquidambar Sweetgums appeared during the mid-late cretaceous and several extinct species are found throughout Asia Europe and North America. The genus was once widespread in europe and asia especially during the miocene. The American Sweetgum is a living fossil itself since fossil specimens dating from the miocene, pliocene and pleistocene were discovered in the eastern united states
Fungi
Neolecta
Animals
Vertebrates
Mammals
Aardvark (Orycteropus afer)
Amami rabbit (Pentalagus furnessi)
Nesolagus (Asian striped rabbits)
Chevrotain (Tragulidae)
Chousingha (Tetracerus quadricornis)
Elephant shrew (Macroscelidea)
Giant panda (Ailuropoda melanoleuca)
Baiji (Lipotes vexillifer) (One living species)
Ganges river dolphin (Platanista gangetica)
Indus river dolphin (Platanista minor)
Hawaiian monk seal (Neomonachus schauinslandi)
Koala (Phascolarctos cinereus)
Laotian rock rat (Laonastes aenigmamus)
Monito del monte (Dromiciops gliroides)
Monotremes (the platypus and echidna)
Mountain beaver (Aplodontia rufa)
Okapi (Okapia johnstoni)
Opossums (Didelphidae)
Clouded leopard (Neofelis nebulousa)
Bush dog (Speothos venaticus)
Maned wolf (Chrysocyon brachyurus)
Red panda (Ailurus fulgens)
Solenodon (Solenodon cubanus and Solenodon paradoxus)
Shrew opossum (Caenolestidae)
Spectacled bear (Tremarctos ornatus)
False killer whale (Pseudorca crassidens)
Pygmy right whale (Caperea marginata)
Pacarana (Dinomys branickii)
Rhinoceroses (Rhinocerotidae)
Tapirs (Tapiridae)
Birds
Pelicans (Pelecanus) – form has been virtually unchanged since the Eocene, and is noted to have been even more conserved across the Cenozoic than that of crocodiles.
Acanthisittidae (New Zealand "wrens") – 2 living species, a few more recently extinct. Distinct lineage of Passeriformes.
Broad-billed sapayoa (Sapayoa aenigma) – One living species. Distinct lineage of Tyranni.
Bearded reedling (Panurus biarmicus) – One living species. Distinct lineage of Passerida or Sylvioidea.
Picathartes (rockfowls)
Coliiformes (mousebirds) – 6 living species in 2 genera. Distinct lineage of Neoaves.
Hoatzin (Ophisthocomus hoazin) – One living species. Distinct lineage of Neoaves.
Magpie goose (Anseranas semipalmata) – One living species. Distinct lineage of Anseriformes.
Sandhill crane (Antigone canadensis) – Oldest living species.
Seriema (Cariamidae) – 2 living species. Distinct lineage of Cariamae.
Tinamiformes (tinamous) 50 living species. Distinct lineage of Palaeognathae.
Reptiles
Crocodilia (crocodiles, gavials, caimans and alligators)
Pig-nosed turtle (Carettochelys insculpta)
Hickatee (Dermatemys mawii)
Snapping turtle (Chelydridae) family
Tuatara (Sphenodon punctatus and Sphenodon guntheri)
Asian forest tortoise (Manouria emys)
Impressed tortoise (Manouria impressa)
Sunbeam snake (Xenopeltis hainanensis and Xenopeltis unicolor)
Leatherback sea turtle (Dermochelys coriacea)
Amphibians
Giant salamanders (Cryptobranchus and Andrias)
Hula painted frog (Latonia nigriventer)
Purple frog (Nasikabatrachus sahyadrensis)
Jawless fish
Hagfish (Myxinidae) family
Lamprey (Petromyzontiformes)
Bony fish
Arowana and arapaima (Osteoglossidae)
Bowfin (Amia calva)
Coelacanth (the lobed-finned Latimeria menadoensis and Latimeria chalumnae)
Gar (Lepisosteidae)
Queensland lungfish (Neoceratodus fosteri)
African lungfish (Protopterus sp.)
Sturgeons and paddlefish (Acipenseriformes)
Bichir (family Polypteridae)
Protanguilla palau
Mudskipper (Oxudercinae)
Sharks
Blind shark (Brachaelurus waddi)
Bullhead shark (Heterodontus sp.)
Cow shark (sixgill sharks and relatives) (Hexanchidae)
Elephant shark (Callorhinchus milii)
Frilled shark (Chlamydoselachus sp.)
Goblin shark (Mitsukurina owstoni)
Gulper shark (Centrophorus sp.)
Invertebrates
Insects
Helorid wasps (1 living genus, 11 extinct genera)
Mantophasmatodea (gladiators; a few living species)
Meropeidae (3 living species, 4 extinct)
Micromalthus debilis (a beetle)
Mymarommatid wasps (10 living species in genus Palaeomymar)
Nevrorthidae (3 species-poor genera)
Nothomyrmecia (known as the 'dinosaur ant')
Notiothauma reedi (a scorpionfly relative)
Orussidae (parasitic wood wasps; about 70 living species in 16 genera)
Peloridiidae (peloridiid bugs; fewer than 30 living species in 13 genera)
Rhinorhipid beetles (1 living species, Triassic origin)
Rotoitid wasps (2 living species, 14 extinct)
Sikhotealinia zhiltzovae (a jurodid beetle)
Syntexis libocedrii (Anaxyelidae cedar wood wasp)
Cyatta abscondita (most recent common relative of Atta and Acromyrmex ant genera)
Crustaceans
Glypheidea (2 living species: Neoglyphea inopinata and Laurentaeglyphea neocaledonica)
Stomatopods (mantis shrimp)
Polychelida (deep sea blind lobster)
Triops cancriformis (also known as tadpole shrimp; a notostracan crustacean)
Molluscs
Nautilina (e.g., Nautilus pompilius)
Neopilina – Monoplacophoran
Slit snail (e.g., Entemnotrochus rumphii)
Vampyroteuthis infernalis – the vampire squid
Pleurocerid snails
Other invertebrates
Crinoids
Springtails
Horseshoe crabs (only 4 living species of the class Xiphosura, family Limulidae)
Lingula anatina (an inarticulate brachiopod)
Liphistiidae (trapdoor spiders)
Onychophorans (velvet worms)
Rhabdopleura (a hemichordate)
Valdiviathyris quenstedti (a craniforman brachiopod)
Paleodictyon nodosum (unknown)
| Biology and health sciences | Basics_4 | Biology |
344287 | https://en.wikipedia.org/wiki/Mercury%20poisoning | Mercury poisoning | Mercury poisoning is a type of metal poisoning due to exposure to mercury. Symptoms depend upon the type, dose, method, and duration of exposure. They may include muscle weakness, poor coordination, numbness in the hands and feet, skin rashes, anxiety, memory problems, trouble speaking, trouble hearing, or trouble seeing. High-level exposure to methylmercury is known as Minamata disease. Methylmercury exposure in children may result in acrodynia (pink disease) in which the skin becomes pink and peels. Long-term complications may include kidney problems and decreased intelligence. The effects of long-term low-dose exposure to methylmercury are unclear.
Forms of mercury exposure include metal, vapor, salt, and organic compound. Most exposure is from eating fish, amalgam-based dental fillings, or exposure at a workplace. In fish, those higher up in the food chain generally have higher levels of mercury, a process known as biomagnification. Less commonly, poisoning may occur as a method of attempted suicide. Human activities that release mercury into the environment include the burning of coal and mining of gold. Tests of the blood, urine, and hair for mercury are available but do not relate well to the amount in the body.
Prevention includes eating a diet low in mercury, removing mercury from medical and other devices, proper disposal of mercury, and not mining further mercury. In those with acute poisoning from inorganic mercury salts, chelation with either dimercaptosuccinic acid (DMSA) or dimercaptopropane sulfonate (DMPS) appears to improve outcomes if given within a few hours of exposure. Chelation for those with long-term exposure is of unclear benefit. In certain communities that survive on fishing, rates of mercury poisoning among children have been as high as 1.7 per 100.
Signs and symptoms
Common symptoms of mercury poisoning are peripheral neuropathy, presenting as paresthesia or itching, burning, pain, or even a sensation that resembles small insects crawling on or under the skin (formication); skin discoloration (pink cheeks, fingertips and toes); swelling; and desquamation (shedding or peeling of skin).
Mercury irreversibly inhibits selenium-dependent enzymes (see below) and may also inactivate S-adenosyl-methionine, which is necessary for catecholamine catabolism by catechol-O-methyl transferase. Due to the body's inability to degrade catecholamines (e.g. adrenaline), a person with mercury poisoning may experience profuse sweating, tachycardia (persistently faster-than-normal heart beat), increased salivation, and hypertension (high blood pressure).
Affected children may show red cheeks, nose and lips, loss of hair, teeth, and nails, transient rashes, hypotonia (muscle weakness), and increased sensitivity to light. Other symptoms may include kidney dysfunction (e.g. Fanconi syndrome) or neuropsychiatric symptoms such as emotional lability, memory impairment, or insomnia.
Thus, the clinical presentation may resemble pheochromocytoma or Kawasaki disease. Desquamation (skin peeling) can occur with severe mercury poisoning acquired by handling elemental mercury.
Causes
Historically, medicines could contain mercury and thus do more harm than good to patients. The popular Victorian medicine calomel contained mercury. In her 1859 autobiography, Scottish seamstress Elizabeth Storie describes her life as a disabled woman due to severe mercury poisoning when a doctor attempted to treat a mild childhood disease with prolonged administration of calomel. In 1862 a soldier in the American civil war, Carleton Burgan, suffered a similar disfigurement when he was treated with calomel for an infection.
Today, consumption of fish containing mercury is by far the most significant source of ingestion-related mercury exposure in humans, although plants and livestock also contain mercury due to bioconcentration of organic mercury from seawater, freshwater, marine and lacustrine sediments, soils, and atmosphere, and due to biomagnification by ingesting other mercury-containing organisms. Exposure to mercury can occur from breathing contaminated air, from eating foods that have acquired mercury residues during processing, from exposure to mercury vapor in mercury amalgam dental restorations, and from improper use or disposal of mercury and mercury-containing objects, for example, after spills of elemental mercury or improper disposal of fluorescent lamps.
All of these, except elemental liquid mercury, produce toxicity or death with less than a gram. Mercury's zero oxidation state (Hg0) exists as vapor or as liquid metal, its mercurous state (Hg+) exists as inorganic salts, and its mercuric state (Hg2+) may form either inorganic salts or organomercury compounds.
Consumption of whale and dolphin meat, as is the practice in Japan, is a source of high levels of mercury poisoning. Tetsuya Endo, a professor at the Health Sciences University of Hokkaido, has tested whale meat purchased in the whaling town of Taiji and found mercury levels more than 20 times the acceptable Japanese standard.
Human-generated sources, such as coal-burning power plants emit about half of atmospheric mercury, with natural sources such as volcanoes responsible for the remainder. A 2021 publication investigating the mercury distribution in European soils found that high mercury concentrations are found close to abandoned mines (such as Almadén (Castilla-La Mancha, Spain), Mt. Amiata (Italy), Idrija (Slovenia) and Rudnany (Slovakia)) and coal-fired power plants. An estimated two-thirds of human-generated mercury comes from stationary combustion, mostly of coal. Other important human-generated sources include gold production, nonferrous metal production, cement production, waste disposal, human crematoria, caustic soda production, pig iron and steel production, mercury production (mostly for batteries), and biomass burning.
Small independent gold-mining operation workers are at higher risk of mercury poisoning because of crude processing methods. Such is the danger for the galamsey in Ghana and similar workers known as orpailleurs in neighboring francophone countries. While no official government estimates of the labor force have been made, observers believe 20,000–50,000 work as galamseys in Ghana, a figure including many women, who work as porters. Similar problems have been reported amongst the gold miners of Indonesia.
Some mercury compounds, especially organomercury compounds, can also be readily absorbed through direct skin contact. Mercury and its compounds are commonly used in chemical laboratories, hospitals, dental clinics, and facilities involved in the production of items such as fluorescent light bulbs, batteries, and explosives.
Many traditional medicines, including ones used in Ayurvedic medicine, and in Traditional Chinese medicine, contain mercury and other heavy metals.
Sources
Organic compounds of mercury tend to be much more toxic than either the elemental form or the salts. These compounds have been implicated in causing brain and liver damage. The most dangerous mercury compound, dimethylmercury, is so toxic that even a few microliters spilled on the skin, or even on a latex glove, can cause death.
Methylmercury and related organomercury compounds
Methylmercury is the major source of organic mercury for all individuals. Due to bioaccumulation, it works its way up through the food web and thus biomagnifies, resulting in high concentrations among populations of some species. Top predatory fish, such as tuna or swordfish, are usually of greater concern than smaller species. The US FDA and the EPA advise women of child-bearing age, nursing mothers, and young children to completely avoid swordfish, shark, king mackerel and tilefish from the Gulf of Mexico, and to limit consumption of albacore ("white") tuna to no more than per week, and of all other fish and shellfish to no more than per week. A 2006 review of the risks and benefits of fish consumption found, for adults, the benefits of one to two servings of fish per week outweigh the risks, even (except for a few fish species) for women of childbearing age, and that avoidance of fish consumption could result in significant excess coronary heart disease deaths and suboptimal neural development in children.
Because the process of mercury-dependent sequestration of selenium is slow, the period between exposure to methylmercury and the appearance of symptoms in adult poisoning cases tends to be extended. The longest recorded latent period is five months after a single exposure, in the Dartmouth case (see History); other latent periods in the range of weeks to months have also been reported. When the first symptom appears, typically paresthesia (a tingling or numbness in the skin), it is followed rapidly by more severe effects, sometimes ending in coma and death. The toxic damage appears to be determined by the peak value of mercury, not the length of the exposure.
Methylmercury exposure during rodent gestation, a developmental period that approximately models human neural development during the first two trimesters of gestation, has long-lasting behavioral consequences that appear in adulthood and, in some cases, may not appear until aging. Prefrontal cortex or dopamine neurotransmission could be especially sensitive to even subtle gestational methylmercury exposure and suggests that public health assessments of methylmercury based on intellectual performance may underestimate the impact of methylmercury in public health.
Ethylmercury is a breakdown product of the antibacteriological agent ethylmercurithiosalicylate, which has been used as a topical antiseptic and a vaccine preservative (further discussed under Thiomersal below). Its characteristics have not been studied as extensively as those of methylmercury. It is cleared from the blood much more rapidly, with a half-life of seven to ten days, and it is metabolized much more quickly than methylmercury. It is presumed not to have methylmercury's ability to cross the blood–brain barrier via a transporter, but instead relies on simple diffusion to enter the brain. Other exposure sources of organic mercury include phenylmercuric acetate and phenylmercuric nitrate. These compounds were used in indoor latex paints for their antimildew properties, but were removed in 1990 because of cases of toxicity.
Inorganic mercury compounds
Mercury occurs as salts such as mercuric chloride (HgCl2) and mercurous chloride (Hg2Cl2), the latter also known as calomel. Because they are more soluble in water, mercuric salts are usually more acutely toxic than mercurous salts. Their higher solubility lets them be more readily absorbed from the gastrointestinal tract. Mercury salts affect primarily the gastrointestinal tract and the kidneys, and can cause severe kidney damage; however, as they cannot cross the blood–brain barrier easily, these salts inflict little neurological damage without continuous or heavy exposure. Mercuric cyanide (Hg(CN)2) is a particularly toxic mercury compound that has been used in murders, as it contains not only mercury but also cyanide, leading to simultaneous cyanide poisoning. The drug n-acetyl penicillamine has been used to treat mercury poisoning with limited success.
Elemental mercury
Quicksilver (liquid metallic mercury) is poorly absorbed by ingestion and skin contact. Its vapor is the most hazardous form. Animal data indicate less than 0.01% of ingested mercury is absorbed through the intact gastrointestinal tract, though it may not be true for individuals with ileus. Cases of systemic toxicity from accidental swallowing are rare, and attempted suicide via intravenous injection does not appear to result in systemic toxicity, though it still causes damage by physically blocking blood vessels both at the site of injection and the lungs. Though not studied quantitatively, the physical properties of liquid elemental mercury limit its absorption through intact skin and in light of its very low absorption rate from the gastrointestinal tract, skin absorption would not be high. Some mercury vapor is absorbed dermally, but uptake by this route is only about 1% of that by inhalation.
In humans, approximately 80% of inhaled mercury vapor is absorbed via the respiratory tract, where it enters the circulatory system and is distributed throughout the body. Chronic exposure by inhalation, even at low concentrations in the range 0.7–42 μg/m3, has been shown in case–control studies to cause effects such as tremors, impaired cognitive skills, and sleep disturbance in workers.
Acute inhalation of high concentrations causes a wide variety of cognitive, personality, sensory, and motor disturbances. The most prominent symptoms include tremors (initially affecting the hands and sometimes spreading to other parts of the body), emotional lability (characterized by irritability, excessive shyness, confidence loss, and nervousness), insomnia, memory loss, neuromuscular changes (weakness, muscle atrophy, muscle twitching), headaches, polyneuropathy (paresthesia, stocking-glove sensory loss, hyperactive tendon reflexes, slowed sensory and motor nerve conduction velocities), and performance deficits in tests of cognitive function.
Mechanism
The toxicity of mercury sources can be expected to depend on its nature, i.e., salts vs. organomercury compounds vs. elemental mercury.
The primary mechanism of mercury toxicity involves its irreversible inhibition of selenoenzymes, such as thioredoxin reductase (IC50 = 9 nM). Although it has many functions, thioredoxin reductase restores vitamins C and E, as well as a number of other important antioxidant molecules, back into their reduced forms, enabling them to counteract oxidative damage. Since the rate of oxygen consumption is particularly high in brain tissues, production of reactive oxygen species (ROS) is accentuated in these vital cells, making them particularly vulnerable to oxidative damage and especially dependent upon the antioxidant protection provided by selenoenzymes. High mercury exposures deplete the amount of cellular selenium available for the biosynthesis of thioredoxin reductase and other selenoenzymes that prevent and reverse oxidative damage, which, if the depletion is severe and long lasting, results in brain cell dysfunctions that can ultimately cause death.
Mercury in its various forms is particularly harmful to fetuses as an environmental toxin in pregnancy, as well as to infants. Women who have been exposed to mercury in substantial excess of dietary selenium intakes during pregnancy are at risk of giving birth to children with serious birth defects, such as those seen in Minamata disease. Mercury exposures in excess of dietary selenium intakes in young children can have severe neurological consequences, preventing nerve sheaths from forming properly.
Exposure to methylmercury causes increased levels of antibodies sent to myelin basic protein (MBP), which is involved in the myelination of neurons, and glial fibrillary acidic protein (GFAP), which is essential to many functions in the central nervous system (CNS). This causes an autoimmmune response against MBP and GFAP and results in the degradation of neural myelin and general decline in function of the CNS.
Diagnosis
Diagnosis of elemental or inorganic mercury poisoning involves determining the history of exposure, physical findings, and an elevated body burden of mercury. Although whole-blood mercury concentrations are typically less than 6 μg/L, diets rich in fish can result in blood mercury concentrations higher than 200 μg/L; it is not that useful to measure these levels for suspected cases of elemental or inorganic poisoning because of mercury's short half-life in the blood. If the exposure is chronic, urine levels can be obtained; 24-hour collections are more reliable than spot collections. It is difficult or impossible to interpret urine samples of people undergoing chelation therapy, as the therapy itself increases mercury levels in the samples.
Diagnosis of organic mercury poisoning differs in that whole-blood or hair analysis is more reliable than urinary mercury levels.
Prevention
Mercury poisoning can be prevented or minimized by eliminating or reducing exposure to mercury and mercury compounds. To that end, many governments and private groups have made efforts to heavily regulate the use of mercury, or to issue advisories about the use of mercury. Most countries have signed the Minamata Convention on Mercury.
The export from the European Union of mercury and some mercury compounds has been prohibited since 15 March 2011. The European Union has banned most uses of mercury. Mercury is allowed for fluorescent light bulbs because of pressure from countries such as Germany, the Netherlands and Hungary, which are connected to the main producers of fluorescent light bulbs: General Electric, Philips and Osram.
The United States Environmental Protection Agency (EPA) issued recommendations in 2004 regarding exposure to mercury in fish and shellfish. The EPA also developed the "Fish Kids" awareness campaign for children and young adults on account of the greater impact of mercury exposure to that population.
Cleaning spilled mercury
Mercury thermometers and mercury light bulbs are not as common as they used to be, and the amount of mercury they contain is unlikely to be a health concern if handled carefully. However, broken items still require careful cleanup, as mercury can be hard to collect and it is easy to accidentally create a much larger exposure problem. If available, powdered sulfur may be applied to the spill, in order to create a solid compound that is more easily removed from surfaces than liquid mercury.
Treatment
Identifying and removing the source of the mercury is crucial. Decontamination requires removal of clothes, washing skin with soap and water, and flushing the eyes with saline solution as needed.
Before the advent of organic chelating agents, salts of iodide were given orally, such as heavily popularized by Louis Melsens and many nineteenth and early twentieth century doctors.
Chelation therapy
Chelation therapy for acute inorganic mercury poisoning, a formerly common method, was done with DMSA, 2,3-dimercapto-1-propanesulfonic acid (DMPS), D-penicillamine (DPCN), or dimercaprol (BAL). Only DMSA is FDA-approved for use in children for treating mercury poisoning. However, several studies found no clear clinical benefit from DMSA treatment for poisoning due to mercury vapor. No chelator for methylmercury or ethylmercury is approved by the FDA; DMSA is the most frequently used for severe methylmercury poisoning, as it is given orally, has fewer side-effects, and has been found to be superior to BAL, DPCN, and DMPS. α-Lipoic acid (ALA) has been shown to be protective against acute mercury poisoning in several mammalian species when it is given soon after exposure; correct dosage is required, as inappropriate dosages increase toxicity. Although it has been hypothesized that frequent low dosages of ALA may have potential as a mercury chelator, studies in rats have been contradictory. Glutathione and N-acetylcysteine (NAC) are recommended by some physicians, but have been shown to increase mercury concentrations in the kidneys and the brain.
Chelation therapy can be hazardous if administered incorrectly. In August 2005, an incorrect form of EDTA (edetate disodium) used for chelation therapy resulted in hypocalcemia, causing cardiac arrest that killed a five-year-old autistic boy.
Other
Experimental animal and epidemiological study findings have confirmed the interaction between selenium and methylmercury. Instead of causing a decline in neurodevelopmental outcomes, epidemiological studies have found that improved nutrient (i.e., omega-3 fatty acids, selenium, iodine, vitamin D) intakes as a result of ocean fish consumption during pregnancy improves maternal and fetal outcomes. For example, increased ocean fish consumption during pregnancy was associated with 4-6 point increases in child IQs.
Prognosis
Some of the toxic effects of mercury are partially or wholly reversible provided specific therapy is able to restore selenium availability to normal before tissue damage from oxidation becomes too extensive. Autopsy findings point to a half-life of inorganic mercury in human brains of 27.4 years. Heavy or prolonged exposure can do irreversible damage, in particular in fetuses, infants, and young children. Young's syndrome is believed to be a long-term consequence of early childhood mercury poisoning.
Mercuric chloride may cause cancer as it has caused increases in several types of tumors in rats and mice, while methyl mercury has caused kidney tumors in male rats. The EPA has classified mercuric chloride and methyl mercury as possible human carcinogens (ATSDR, EPA)
Detection in biological fluids
Mercury may be measured in blood or urine to confirm a diagnosis of poisoning in hospitalized people or to assist in the forensic investigation in a case of fatal over dosage. Some analytical techniques are capable of distinguishing organic from inorganic forms of the metal. The concentrations in both fluids tend to reach high levels early after exposure to inorganic forms, while lower but very persistent levels are observed following exposure to elemental or organic mercury. Chelation therapy can cause a transient elevation of urine mercury levels.
History
Neolithic artists using cinnabar show signs of mercury poisoning.
Several Chinese emperors and other Chinese nobles are known or suspected to have died or been sickened by mercury poisoning after alchemists administered them "elixirs" to promote health, longevity, or immortality that contained either elemental mercury or (more commonly) cinnabar. Among the most prominent examples:
The first emperor of unified China, Qin Shi Huang, it is reported, died in 210 BC of ingesting mercury pills that were intended to give him eternal life.
Emperor Xuānzong of Tang, one of the emperors of the late Tang dynasty of China, was prescribed "cinnabar that had been treated and subdued by fire" to achieve immortality. Concerns that the prescription was having ill effects on the emperor's health and sanity were waved off by the imperial alchemists, who cited medical texts listing a number of the emperor's conditions (including itching, formication, swelling, and muscle weakness), today recognized as signs and symptoms of mercury poisoning, as evidence that the elixir was effectively treating the emperor's latent ailments. Xuānzong became irritable and paranoid, and he seems to have ultimately died in 859 from the poisoning.
In his Natural History, Pliny the Elder writes that "it is a fact generally admitted that [cinnabar] is a poison" and warns against using it in medicine, also noting that workers polishing it "tie on their face loose masks of bladder-skin, to prevent their inhaling the dust in breathing", one of the earliest mentions of PPE.
Carl Scheele, a significant 18th century Swedish pioneer of chemical research, died from mercury poisoning arising from his work, at the relatively early age of 43.
The phrase mad as a hatter is likely a reference to mercury poisoning among milliners (so-called "mad hatter disease"), as mercury-based compounds were once used in the manufacture of felt hats in the 18th and 19th century. (The Mad Hatter character of Alice in Wonderland was, it is presumed, inspired by an eccentric furniture dealer named Theophilus Carter. Carter was not a victim of mad hatter disease although Lewis Carroll would have been familiar with the phenomenon of dementia that occurred among hatters.)
In 1810, two British ships, HMS Triumph and , salvaged a large load of elemental mercury from a wrecked Spanish vessel near Cadiz, Spain. The bladders containing the mercury soon ruptured. The element spread about the ships in liquid and vapor forms. The sailors presented with neurologic compromises: tremor, paralysis, and excessive salivation as well as tooth loss, skin problems, and pulmonary complaints. In 1823 William Burnett, M.D. published a report on the effects of mercurial vapor. Triumph surgeon, Henry Plowman, had concluded that the ailments had arisen from inhaling the mercurialized atmosphere. His treatment was to order the lower deck gun ports to be opened, when it was safe to do so; sleeping on the orlop was forbidden; and no men slept in the lower deck if they were at all symptomatic. Windsails were set to channel fresh air into the lower decks day and night.
Historically, gold-mercury amalgam was widely used in gilding, applied to the object and then heated to vaporize the mercury and deposit the gold, leading to numerous casualties among the workers. It is estimated that during the construction of Saint Isaac's Cathedral alone, 60 men died from the gilding of the main dome.
For years, including the early part of his presidency, Abraham Lincoln took a common medicine of his time called "blue mass", which contained significant amounts of mercury.
On September 5, 1920, silent movie actress Olive Thomas ingested mercury capsules dissolved in an alcoholic solution at the Hotel Ritz in Paris. There is still controversy over whether it was suicide, or whether she consumed the external preparation by mistake. Her husband, Jack Pickford (the brother of Mary Pickford), had syphilis, and the mercury was used as a treatment of the venereal disease at the time. She died a few days later at the American Hospital in Neuilly.
An early scientific study of mercury poisoning was in 1923–1926 by the German inorganic chemist, Alfred Stock, who himself became poisoned, together with his colleagues, by breathing mercury vapor that was being released by his laboratory equipment—diffusion pumps, float valves, and manometers—all of which contained mercury, and also from mercury that had been accidentally spilt and remained in cracks in the linoleum floor covering. He published a number of papers on mercury poisoning, founded a committee in Berlin to study cases of possible mercury poisoning, and introduced the term micromercurialism.
The term Hunter-Russell syndrome derives from a study of mercury poisoning among workers in a seed-packaging factory in Norwich, England in the late 1930s who breathed methylmercury that was being used as a seed disinfectant and pesticide.
Outbreaks of methylmercury poisoning occurred in several places in Japan during the 1950s due to industrial discharges of mercury into rivers and coastal waters. The best-known instances were in Minamata and Niigata. In Minamata alone, more than 600 people died due to what became known as Minamata disease. More than 21,000 people filed claims with the Japanese government, of which almost 3000 became certified as having the disease. In 22 documented cases, pregnant women who consumed contaminated fish showed mild or no symptoms but gave birth to infants with severe developmental disabilities.
Mercury poisoning of generations of Grassy Narrows and Whitedog native people in Ontario, Canada who were exposed to high levels of mercury by consuming mercury-contaminated fish when Dryden Chemical Company discharged over of mercury directly into the Wabigoon–English River system and continued with mercury air pollution until 1975.
Widespread mercury poisoning occurred in rural Iraq in 1971–1972, when grain treated with a methylmercury-based fungicide that was intended for planting only was used by the rural population to make bread, causing at least 6530 cases of mercury poisoning and at least 459 deaths (see Basra poison grain disaster).
On August 14, 1996, Karen Wetterhahn, a chemistry professor working at Dartmouth College, spilled a small amount of dimethylmercury on her latex glove. She began experiencing the symptoms of mercury poisoning five months later and, despite aggressive chelation therapy, died a few months later from a mercury induced neurodegenerative disease
In April 2000, Alan Chmurny attempted to kill a former employee, Marta Bradley, by pouring mercury into the ventilation system of her car.
On March 19, 2008, Tony Winnett, 55, inhaled mercury vapors while trying to extract gold from computer parts (by using liquid mercury to separate gold from the rest of the alloy), and died ten days later. His Oklahoma residence became so contaminated that it had to be gutted.
In December 2008, actor Jeremy Piven was diagnosed with mercury poisoning possibly resulting from eating sushi twice a day for twenty years or from taking herbal remedies.
In India, a study by Centre for Science and Environment and Indian Institute of Toxicology Research has found that in the country's energy capital Singrauli, mercury is slowly entering people's homes, food, water and even blood.
The Minamata Convention on Mercury in 2016 announced that the signing of the "international treaty designed to protect human health and the environment from anthropogenic releases and emission of mercury and mercury compounds" on April 22, 2016—Earth Day. It was the sixtieth anniversary of the discovery of the disease.
In August 2024, chess player Amina Abakarova allegedly attempted to poison her rival, Umayganat Osmanova, by coating chess pieces in mercury from a thermometer.
Infantile acrodynia
Infantile acrodynia (also known as "calomel disease", "erythredemic polyneuropathy", and "pink disease") is a type of mercury poisoning in children characterized by pain and pink discoloration of the hands and feet. The word is derived from the Greek, where άκρο means end or extremity, and οδυνη means pain. Acrodynia resulted primarily from calomel in teething powders and decreased greatly after calomel was excluded from most teething powders in 1954.
Acrodynia is difficult to diagnose; "it is most often postulated that the etiology of this syndrome is an idiosyncratic hypersensitivity reaction to mercury because of the lack of correlation with mercury levels, many of the symptoms resemble recognized mercury poisoning."
Medicine
Mercury was once prescribed as a purgative.
Many mercury-containing compounds were once used in medicines. These include calomel (mercurous chloride), and mercuric chloride.
Thiomersal
In 1999, the Centers for Disease Control (CDC) and the American Academy of Pediatrics (AAP) asked vaccine makers to remove the organomercury compound thiomersal (spelled "thimerosal" in the US) from vaccines as quickly as possible, and thiomersal has been phased out of US and European vaccines, except for some preparations of influenza vaccine. The CDC and the AAP followed the precautionary principle, which assumes that there is no harm in exercising caution even if it later turns out to be unwarranted, but their 1999 action sparked confusion and controversy that thiomersal was a cause of autism.
Since 2000, the thiomersal in child vaccines has been alleged to contribute to autism, and thousands of parents in the United States have pursued legal compensation from a federal fund. A 2004 Institute of Medicine (IOM) committee favored rejecting any causal relationship between thiomersal-containing vaccines and autism. Autism incidence rates increased steadily even after thiomersal was removed from childhood vaccines. Currently there is no accepted scientific evidence that exposure to thiomersal is a factor in causing autism.
Dental amalgam toxicity
Dental amalgam is a possible cause of low-level mercury poisoning due to its use in dental fillings. Discussion on the topic includes debates on whether amalgam should be used, with critics arguing that its toxic effects make it unsafe.
Cosmetics
Some skin whitening products contain the toxic mercury(II) chloride as the active ingredient. When applied, the chemical readily absorbs through the skin into the bloodstream. The use of mercury in cosmetics is illegal in the United States. However, cosmetics containing mercury are often illegally imported. Following a certified case of mercury poisoning resulting from the use of an imported skin whitening product, the United States Food and Drug Administration warned against the use of such products. Symptoms of mercury poisoning have resulted from the use of various mercury-containing cosmetic products. The use of skin whitening products is especially popular amongst Asian women. In Hong Kong in 2002, two products were discovered to contain between 9,000 and 60,000 times the recommended dose.
Fluorescent lamps
Fluorescent lamps contain mercury, which is released when bulbs break. Mercury in bulbs is typically present as either elemental mercury liquid, vapor, or both, since the liquid evaporates at ambient temperature. When broken indoors, bulbs may emit sufficient mercury vapor to present health concerns, and the U.S. Environmental Protection Agency recommends evacuating and airing out a room for at least 15 minutes after breaking a fluorescent light bulb. Breakage of multiple bulbs presents a greater concern. A 1987 report described a 23-month-old toddler who had anorexia, weight loss, irritability, profuse sweating, and peeling and redness of fingers and toes. This case of acrodynia was traced to exposure of mercury from a carton of 8-foot fluorescent light bulbs that had broken in a potting shed adjacent to the main nursery. The glass was cleaned up and discarded, but the child often used the area to play in.
Assassination attempts
Mercury has, allegedly, been used at various times to assassinate people. In 2008, Russian lawyer Karinna Moskalenko claimed to have been poisoned by mercury left in her car, while in 2010 journalists Viktor Kalashnikov and Marina Kalashnikova accused Russia's FSB of trying to poison them.
In 2011, German Christoph Bulwin was poisoned with a mercury compound from a syringe attached to an umbrella.
| Biology and health sciences | Types | Health |
11946500 | https://en.wikipedia.org/wiki/HD%20189733%20b | HD 189733 b | HD 189733 b is an exoplanet in the constellation of Vulpecula approximately away from the Solar System. Astronomers in France discovered the planet orbiting the star HD 189733 on October 5, 2005, by observing its transit across the star's face. With a mass 11.2% higher than that of Jupiter and a radius 11.4% greater, HD 189733 b orbits its host star once every 2.2 days at an orbital speed of , making it a hot Jupiter with poor prospects for extraterrestrial life.
The closest transiting hot Jupiter to Earth, HD 189733 b has been the subject of close atmospheric observation. Scientists have studied it with high- and low-resolution instruments, both from the ground and from space. Researchers have found that the planet's weather includes raining molten glass. HD 189733 b was also the first exoplanet to have its thermal map constructed, possibly to be detected through polarimetry, its overall color determined (deep blue), its transit viewed in the X-ray spectrum, and to have carbon dioxide confirmed as being present in its atmosphere.
In July 2014, NASA announced the discovery of very dry atmospheres on three exoplanets that orbited Sun-like stars: HD 189733 b, HD 209458 b, and WASP-12b.
Detection and discovery
Transit and Doppler spectroscopy
On October 6, 2005, a team of astronomers announced the discovery of transiting planet HD 189733 b. The planet was then detected using Doppler spectroscopy. Real-time radial velocity measurements detected the Rossiter–McLaughlin effect caused by the planet passing in front of its star before photometric measurements confirmed that the planet was transiting. In 2006, a team led by Drake Deming announced detection of strong infrared thermal emission from the transiting exoplanet planet HD 189733 b, by measuring the flux decrement (decrease of total light) during its prominent secondary eclipse (when the planet passes behind the star).
The mass of the planet is estimated to be 16% larger than Jupiter's, with the planet completing an orbit around its host star every 2.2 days and an orbital speed of .
Infrared spectrum
On February 21, 2007, NASA released news that the Spitzer Space Telescope had measured detailed spectra from both HD 189733 b and HD 209458 b. The release came simultaneously with the public release of a new issue of Nature containing the first publication on the spectroscopic observation of the other exoplanet, HD 209458 b. A paper was submitted and published by the Astrophysical Journal Letters. The spectroscopic observations of HD 189733 b were led by Carl Grillmair of NASA's Spitzer Science Center.
Visible color
In 2008, a team of astrophysicists appeared to have detected and monitored the planet's visible light using polarimetry, which would have been the first such success. This result seemed to be confirmed and refined by the same team in 2011. They found that the planet albedo is significantly larger in blue light than in the red, most probably due to Rayleigh scattering and molecular absorption in the red. The blue color of the planet was subsequently confirmed in 2013, which would have made HD 189733 the first planet to have its overall color determined by two different techniques. The measurements in polarized light have since been disputed by two separate teams using more sensitive polarimeters, with upper limits of the polarimetric signal provided therein.
The rich cobalt blue colour of HD 189733 b may be the result of Rayleigh scattering. In mid January 2008, spectral observation during the planet's transit using that model found that if molecular hydrogen exists, it would have an atmospheric pressure of 410 ± 30 mbar of 0.1564 solar radii. The Mie approximation model also found that there is a possible condensate in its atmosphere, magnesium silicate (MgSiO3) with a particle size of approximately 10−2 to 10−1 μm. Using both models, the planet's temperature would be between 1340 and 1540 K. The Rayleigh effect is confirmed in other models, and by the apparent lack of a cooler, shaded stratosphere below its outer atmosphere. In the visible region of the spectrum, thanks to their high absorption cross sections, atomic sodium and potassium can be investigated. For example, using high-resolution UVES spectrograph on the Very Large Telescope, sodium has been detected on this atmosphere and further physical characteristics of the atmosphere such as temperature has been investigated.
X-ray spectrum
In July 2013, NASA reported the first observations of planet transit studied in the X-ray spectrum. It was found that the planet's atmosphere blocks three times more X-rays than visible light.
Evaporation
In March 2010, transit observations using HI Lyman-alpha found that this planet is evaporating at a rate of 1-100 gigagrams per second. This indication was found by detecting the extended exosphere of atomic hydrogen. HD 189733 b is the second planet after HD 209458 b for which atmospheric evaporation has been detected.
Physical characteristics
This planet exhibits one of the largest photometric transit depth (amount of the parent star's light blocked) of extrasolar planets so far observed, approximately 3%. The apparent longitude of ascending node of its orbit is 16 degrees +/- 8 away from the north–south in our sky. It and HD 209458 b were the first two planets to be directly spectroscopically observed. The parent stars of these two planets are the brightest transiting-planet host stars, so these planets will continue to receive the most attention from astronomers. Like most hot Jupiters, this planet is thought to be tidally locked to its parent star, meaning it has a permanent day and night.
The planet is not oblate, and has neither satellites with greater than 0.8 the radius of Earth nor a ring system like that of Saturn.
The international team under the direction of Svetlana Berdyugina of Zurich University of Technology, using the Swedish 60-centimeter telescope KVA, which is located in Spain, was able to directly see the polarized light reflected from the planet. The polarization indicates that the scattering atmosphere is considerably larger (> 30%) than the opaque body of the planet seen during transits.
The atmosphere was at first predicted "pL class", lacking a temperature-inversion stratosphere; like L dwarfs which lack titanium and vanadium oxides. Follow-up measurements, tested against a stratospheric model, yielded inconclusive results. Atmospheric condensates form a haze above the surface as viewed in the infrared. A sunset viewed from that surface would be red. Sodium and potassium signals were predicted by Tinetti 2007. First obscured by the haze of condensates, sodium was eventually observed at three times the concentration of HD 209458 b's sodium layer. The potassium was also detected in 2020, although in significantly smaller concentrations. HD 189733 is also the first extrasolar planet confirmed to have carbon dioxide in its atmosphere. In 2024, hydrogen sulfide was detected in HD 189733 b's atmosphere.
Map of the planet
In 2007, the Spitzer Space Telescope was used to map the planet's temperature emissions. The planet and star system was observed for 33 consecutive hours, starting when only the night side of the planet was in view. Over the course of one-half of the planet's orbit, more and more of the dayside came into view. A temperature range of 973 ± 33 K to 1,212 ± 11 K was discovered, indicating that the absorbed energy from the parent star is distributed fairly evenly through the planet's atmosphere. The region of peak temperature was offset 30 degrees east of the substellar point, as predicted by theoretical models of hot Jupiters taking into account a parameterized day to night redistribution mechanism.
Scientists at the University of Warwick determined that HD 189733 b has winds of up to blowing from the day side to the night side. NASA released a brightness map of the surface temperature of HD 189733 b; it is the first map ever published of an extra-solar planet.
Water vapor, oxygen, and organic compounds
On July 11, 2007, a team led by Giovanna Tinetti published the results of their observations using the Spitzer Space Telescope concluding there is solid evidence for significant amounts of water vapor in the planet's atmosphere. Follow-up observations made using the Hubble Space Telescope confirm the presence of water vapor, neutral oxygen and also the organic compound methane. Later, Very Large Telescope observations also detected the presence of carbon monoxide on the day side of the planet. It is currently unknown how the methane originated as the planet's high temperature should cause the water and methane to react, replacing the atmosphere with carbon monoxide. Nonetheless, the presence of roughly 0.004% of water vapour fraction by volume in atmosphere of HD 189733 b was confirmed with high-resolution emission spectra taken in 2021.
Evolution
While transiting the system also clearly exhibits the Rossiter–McLaughlin effect, shifting in photospheric spectral lines caused by the planet occulting a part of the rotating stellar surface. Due to its high mass and close orbit, the parent star has a very large semi-amplitude (K), the "wobble" in the star's radial velocity, of 205 m/s.
The Rossiter–McLaughlin effect allows the measurement of the angle between the planet's orbital plane and the equatorial plane of the star. These are well aligned, misalignment equal to -0.5°. By analogy with HD 149026 b, the formation of the planet was peaceful and probably involved interactions with the protoplanetary disc. A much larger angle would have suggested a violent interplay with other protoplanets.
Star-planet interaction controversy
In 2008, a team of astronomers first described how as the exoplanet orbiting HD 189733 A reaches a certain place in its orbit, it causes increased stellar flaring. In 2010, a different team found that every time they observe the exoplanet at a certain position in its orbit, they also detected X-ray flares. Theoretical research since 2000 suggested that an exoplanet very near to the star that it orbits may cause increased flaring due to the interaction of their magnetic fields, or because of tidal forces. In 2019, astronomers analyzed data from Arecibo Observatory, MOST, and the Automated Photoelectric Telescope, in addition to historical observations of the star at radio, optical, ultraviolet, and X-ray wavelengths to examine these claims. They found that the previous claims were exaggerated and the host star failed to display many of the brightness and spectral characteristics associated with stellar flaring and solar active regions, including sunspots. Their statistical analysis also found that many stellar flares are seen regardless of the position of the exoplanet, therefore debunking the earlier claims. The magnetic fields of the host star and exoplanet do not interact, and this system is no longer believed to have a "star-planet interaction." Some researchers had also suggested that HD 189733 accretes, or pulls, gas from its orbiting exoplanet at a rate similar to those found around young protostars in T Tauri Star systems. Later analysis demonstrated that very little, if any, gas was accreted from the "hot Jupiter" companion.
Possible exomoons
Some studies have proposed candidate exomoons around HD 189733 b. A 2014 study proposed a moon based on studying periodic increases and decreases in light given off from HD 189733 b. This moon would be outside of the planet's Hill sphere, making its existence implausible. Two studies by the same team in 2019 and 2020 proposed exo-Io candidates around a number of hot Jupiters, including HD 189733 b and WASP-49b, based on detected sodium and potassium, consistent with evaporating exomoons and/or their corresponding gas torus. A follow-up study in 2022 did not find evidence for an exomoon around HD 189733 b.
| Physical sciences | Notable exoplanets | Astronomy |
9491546 | https://en.wikipedia.org/wiki/Slickenside | Slickenside | In geology, a slickenside is a smoothly polished surface caused by frictional movement between rocks along a fault. This surface is typically striated with linear features, called slickenlines, in the direction of movement.
Geometry of slickensides
A slickenside can occur as a single surface at a fault between two hard surfaces. Alternatively, the gouge between the fault surfaces may contain many anastamosing slip surfaces that host slickensides. These slip surfaces are on the order of 100 micrometers thick, and the size of the grains that constitute the surface are ultra-fine (0.01-1 micrometers in diameter). These grains are unlike typical grains of fault rock in that they have irregular grain boundaries and few crystal lattice defects (termed dislocations).
Slickensides have conspicuous shapes that can be used to determine the direction of movement along the fault. Straight slickenlines indicate linear-translational fault motion. They are parallel to the direction of fault motion and serve as a kinematic indicator. Curved slickenlines have recently been studied for their potential to preserve the direction of earthquake rupture propagation.
Surface roughness
Slickenside formation results in unique roughness on a slip surface. Fault surface roughness (or topography) is characterized by the aspect ratio of asperity height to scale of observation, and this roughness is a key parameter in the study of fault slip. In general, a fault surface appears rougher at smaller scales (i.e. rough and bumpy at approximately millimetre scales and smaller, and increasingly smooth with larger fields of view). This smoothing with larger observation scales is more pronounced in the slip-parallel direction than the slip-perpendicular direction and is commonly a result of slickenside formation.
Mechanisms to create slickensides
The unique geometry of a slickenside can be created in a variety of ways, but the precise mechanisms that create them is not well understood. The grinding between two rocks produces granular material, and there is a change in the behaviour of wear material when the particle size is reduced to nanometers. When the particle size is reduced so dramatically that the surface becomes shiny, it can be characterized as a fault mirror.
A fault mirror may also be the result of fluid being present at the fault surface during slip. Once slip has stopped, this fluid solidifies as a silica gel, which appears shiny and hosts slickenlines.
Asperity plowing
An asperity on a fault surface is a bump or point with higher relief than the area around it. The asperity, when pressed into the opposing rock surface and then moved, digs into the opposing rock, forming troughs, grooves, and scratches. Asperity plowing is thus a result of permanent deformation in the brittle regime at a small scale.
Debris streaking
When an asperity plows into the opposing rock, it wears itself and the opposing rock down and produces fine debris. This debris, or wear product, accumulates both in front of and behind the asperity in a long, elongated shape. If the asperity is relatively hard, the debris will accumulate in front of the asperity. If the asperity is relatively soft, the debris will trail behind. This debris hardens over time and is preserved as a form of slickenline.
Erosional sheltering
Some rocks may contain particles that are harder than the rest of the rock. When these rocks are worn, the harder particles will resist wear more than the softer rock, the rock on the lee side of the hard particle will be protected from wear. This creates a tail that starts abruptly as a crag where the hard particle was located and is elongated parallel to the direction of movement down-slip from the particle.
Fibre growth
The fault plane may be coated by mineral fibres that grew in during the fault movement, known as slickenfibres. Due to irregularities in the fault plane, exposed slickenfibres typically have a stepped appearance that can be used to determine the sense of movement across the fault.
Slickenfibres are secondary minerals that make up the slickensides rather than the rock itself. Slickenfibres form in areas where the rock slowly creep past one another rather than sliding suddenly as a result of an earthquake. Unlike slickenlines, which give two possibilities for slip direction, slickenfibres preserve the true slip direction.
Implications
Slickensides provide useful insight into earthquake processes. Calcite slickenfibres have recently been used to constrain the depth of aseismic creep in the Zagros Mountains as well as the orientation of stress acting on the fault. It has also been suggested that when multiple slickenfibre or slickensteps orientations are present, it can indicate that the ongoing shear is not strain softening so slip does not have a constant direction.
In addition to the direction of slip, slickenlines have also been used to constrain the timing of fault slip. They also preserve any complexity in the geometry of the earthquake rupture.
Other types of slickensides
Slickensides in soils
In pedology, the study of soils in their natural environments, a slickenside is a surface of the cracks produced in soils containing a high proportion of swelling clays. Slickensides are a type of cutan.
In the Australian Soil Classification, slickensides, along with lenticular structural aggregates, are an indicator of a vertisol.
Slickensides on the Moon
On the Moon, a boulder with slickensides, discovered in a debris-strewn small crater at Station 9 near Rima Hadley, was photographed during a moonwalk by the crew of Apollo 15.
Gallery
| Physical sciences | Geologic features | Earth science |
2229292 | https://en.wikipedia.org/wiki/Stirling%20numbers%20of%20the%20second%20kind | Stirling numbers of the second kind | In mathematics, particularly in combinatorics, a Stirling number of the second kind (or Stirling partition number) is the number of ways to partition a set of n objects into k non-empty subsets and is denoted by or . Stirling numbers of the second kind occur in the field of mathematics called combinatorics and the study of partitions. They are named after James Stirling.
The Stirling numbers of the first and second kind can be understood as inverses of one another when viewed as triangular matrices. This article is devoted to specifics of Stirling numbers of the second kind. Identities linking the two kinds appear in the article on Stirling numbers.
Definition
The Stirling numbers of the second kind, written or or with other notations, count the number of ways to partition a set of labelled objects into nonempty unlabelled subsets. Equivalently, they count the number of different equivalence relations with precisely equivalence classes that can be defined on an element set. In fact, there is a bijection between the set of partitions and the set of equivalence relations on a given set. Obviously,
for n ≥ 0, and for n ≥ 1,
as the only way to partition an n-element set into n parts is to put each element of the set into its own part, and the only way to partition a nonempty set into one part is to put all of the elements in the same part. Unlike Stirling numbers of the first kind, they can be calculated using a one-sum formula:
The Stirling numbers of the first kind may be characterized as the numbers that arise when one expresses powers of an indeterminate x in terms of the falling factorials
(In particular, (x)0 = 1 because it is an empty product.)
Stirling numbers of the second kind satisfy the relation
Notation
Various notations have been used for Stirling numbers of the second kind. The brace notation was used by Imanuel Marx and Antonio Salmeri in 1962 for variants of these numbers.<ref>Antonio Salmeri, Introduzione alla teoria dei coefficienti fattoriali, Giornale di Matematiche di Battaglini 90 (1962), pp. 44–54.</ref> This led Knuth to use it, as shown here, in the first volume of The Art of Computer Programming (1968).Donald E. Knuth, Fundamental Algorithms, Reading, Mass.: Addison–Wesley, 1968. According to the third edition of The Art of Computer Programming, this notation was also used earlier by Jovan Karamata in 1935.Jovan Karamata, Théorèmes sur la sommabilité exponentielle et d'autres sommabilités s'y rattachant, Mathematica (Cluj) 9 (1935), pp, 164–178. The notation S(n, k) was used by Richard Stanley in his book Enumerative Combinatorics and also, much earlier, by many other writers.
The notations used on this page for Stirling numbers are not universal, and may conflict with notations in other sources.
Relation to Bell numbers
Since the Stirling number counts set partitions of an n-element set into k parts, the sum
over all values of k is the total number of partitions of a set with n members. This number is known as the nth Bell number.
Analogously, the ordered Bell numbers can be computed from the Stirling numbers of the second kind via
Table of values
Below is a triangular array of values for the Stirling numbers of the second kind :
As with the binomial coefficients, this table could be extended to , but the entries would all be 0.
Properties
Recurrence relation
Stirling numbers of the second kind obey the recurrence relation
with initial conditions
For instance, the number 25 in column k = 3 and row n = 5 is given by 25 = 7 + (3×6), where 7 is the number above and to the left of 25, 6 is the number above 25 and 3 is the column containing the 6.
To prove this recurrence, observe that a partition of the objects into k nonempty subsets either contains the -th object as a singleton or it does not. The number of ways that the singleton is one of the subsets is given by
since we must partition the remaining objects into the available subsets. In the other case the -th object belongs to a subset containing other objects. The number of ways is given by
since we partition all objects other than the -th into k subsets, and then we are left with k choices for inserting object . Summing these two values gives the desired result.
Another recurrence relation is given by
which follows from evaluating at .
Simple identities
Some simple identities include
This is because dividing n elements into sets necessarily means dividing it into one set of size 2 and sets of size 1. Therefore we need only pick those two elements;
and
To see this, first note that there are 2 ordered pairs of complementary subsets A and B. In one case, A is empty, and in another B is empty, so ordered pairs of subsets remain. Finally, since we want unordered pairs rather than ordered pairs we divide this last number by 2, giving the result above.
Another explicit expansion of the recurrence-relation gives identities in the spirit of the above example.
Identities
The table in section 6.1 of Concrete Mathematics provides a plethora of generalized forms of finite sums involving the Stirling numbers. Several particular finite sums relevant to this article include
Explicit formula
The Stirling numbers of the second kind are given by the explicit formula:
This can be derived by using inclusion-exclusion to count the surjections from n to k and using the fact that the number of such surjections is .
Additionally, this formula is a special case of the kth forward difference of the monomial evaluated at x = 0:
Because the Bernoulli polynomials may be written in terms of these forward differences, one immediately obtains a relation in the Bernoulli numbers:
The evaluation of incomplete exponential Bell polynomial Bn,k(x1,x2,...) on the sequence of ones equals a Stirling number of the second kind:
Another explicit formula given in the NIST Handbook of Mathematical Functions is
Parity
The parity of a Stirling number of the second kind is same as the parity of a related binomial coefficient:
where
This relation is specified by mapping n and k coordinates onto the Sierpiński triangle.
More directly, let two sets contain positions of 1's in binary representations of results of respective expressions:
One can mimic a bitwise AND operation by intersecting these two sets:
to obtain the parity of a Stirling number of the second kind in O(1) time. In pseudocode:
where is the Iverson bracket.
The parity of a central Stirling number of the second kind is odd if and only if is a fibbinary number, a number whose binary representation has no two consecutive 1s.
Generating functions
For a fixed integer n, the ordinary generating function for Stirling numbers of the second kind is given by
where are Touchard polynomials. If one sums the Stirling numbers against the falling factorial instead, one can show the following identities, among others:
and
which has special case
For a fixed integer k, the Stirling numbers of the second kind have rational ordinary generating function
and have an exponential generating function given by
A mixed bivariate generating function for the Stirling numbers of the second kind is
Lower and upper bounds
If and , then
Asymptotic approximation
For fixed value of the asymptotic value of the Stirling numbers of the second kind as is given by
If (where o denotes the little o notation) then
A uniformly valid approximation also exists: for all such that , one has
where , and is the unique solution to . Relative error is bounded by about .
Unimodality
For fixed , is unimodal, that is, the sequence increases and then decreases. The maximum is attained for at most two consecutive values of k. That is, there is an integer such that
Looking at the table of values above, the first few values for are
When is large
and the maximum value of the Stirling number can be approximated with
Applications
Moments of the Poisson distribution
If X is a random variable with a Poisson distribution with expected value λ, then its n-th moment is
In particular, the nth moment of the Poisson distribution with expected value 1 is precisely the number of partitions of a set of size n, i.e., it is the nth Bell number (this fact is Dobiński's formula).
Moments of fixed points of random permutations
Let the random variable X be the number of fixed points of a uniformly distributed random permutation of a finite set of size m. Then the nth moment of X is
Note: The upper bound of summation is m, not n.
In other words, the nth moment of this probability distribution is the number of partitions of a set of size n into no more than m parts.
This is proved in the article on random permutation statistics, although the notation is a bit different.
Rhyming schemes
The Stirling numbers of the second kind can represent the total number of rhyme schemes for a poem of n lines. gives the number of possible rhyming schemes for n lines using k unique rhyming syllables. As an example, for a poem of 3 lines, there is 1 rhyme scheme using just one rhyme (aaa), 3 rhyme schemes using two rhymes (aab, aba, abb), and 1 rhyme scheme using three rhymes (abc).
Variants
r-Stirling numbers of the second kind
The r-Stirling number of the second kind counts the number of partitions of a set of n objects into k non-empty disjoint subsets, such that the first r elements are in distinct subsets. These numbers satisfy the recurrence relation
Some combinatorial identities and a connection between these numbers and context-free grammars can be found in
Associated Stirling numbers of the second kind
An r-associated Stirling number of the second kind is the number of ways to partition a set of n objects into k subsets, with each subset containing at least r elements. It is denoted by and obeys the recurrence relation
The 2-associated numbers appear elsewhere as "Ward numbers" and as the magnitudes of the coefficients of Mahler polynomials.
Reduced Stirling numbers of the second kind
Denote the n objects to partition by the integers 1, 2, ..., n. Define the reduced Stirling numbers of the second kind, denoted , to be the number of ways to partition the integers 1, 2, ..., n into k nonempty subsets such that all elements in each subset have pairwise distance at least d. That is, for any integers i and j in a given subset, it is required that . It has been shown that these numbers satisfy
(hence the name "reduced"). Observe (both by definition and by the reduction formula), that , the familiar Stirling numbers of the second kind.
| Mathematics | Combinatorics | null |
2229296 | https://en.wikipedia.org/wiki/Stirling%20numbers%20of%20the%20first%20kind | Stirling numbers of the first kind | In mathematics, especially in combinatorics, Stirling numbers of the first kind arise in the study of permutations. In particular, the unsigned Stirling numbers of the first kind count permutations according to their number of cycles (counting fixed points as cycles of length one).
The Stirling numbers of the first and second kind can be understood as inverses of one another when viewed as triangular matrices. This article is devoted to specifics of Stirling numbers of the first kind. Identities linking the two kinds appear in the article on Stirling numbers.
Definitions
Definition by algebra
Signed Stirling numbers of the first kind are the coefficients in the expansion of the falling factorial
into powers of the variable :
For example, , leading to the values , , and .
The unsigned Stirling numbers may also be defined algebraically as the coefficients of the rising factorial:
.
The notations used on this page for Stirling numbers are not universal, and may conflict with notations in other sources; the square bracket notation is also common notation for the Gaussian coefficients.
Definition by permutation
Subsequently, it was discovered that the absolute values of these numbers are equal to the number of permutations of certain kinds. These absolute values, which are known as unsigned Stirling numbers of the first kind, are often denoted or . They may be defined directly to be the number of permutations of elements with disjoint cycles.
For example, of the permutations of three elements, there is one permutation with three cycles (the identity permutation, given in one-line notation by or in cycle notation by ), three permutations with two cycles (, , and ) and two permutations with one cycle ( and ). Thus , and . These can be seen to agree with the previous algebraic calculations of for .
For another example, the image at right shows that : the symmetric group on 4 objects has 3 permutations of the form
(having 2 orbits, each of size 2),
and 8 permutations of the form
(having 1 orbit of size 3 and 1 orbit of size 1).
These numbers can be calculated by considering the orbits as conjugacy classes. Alfréd Rényi observed that the unsigned Stirling number of the first kind also counts the number of permutations of size with left-to-right maxima.
Signs
The signs of the signed Stirling numbers of the first kind depend only on the parity of .
Recurrence relation
The unsigned Stirling numbers of the first kind follow the recurrence relation
for , with the boundary conditions
for .
It follows immediately that the signed Stirling numbers of the first kind satisfy the recurrence
.
Table of values
Below is a triangular array of unsigned values for the Stirling numbers of the first kind, similar in form to Pascal's triangle. These values are easy to generate using the recurrence relation in the previous section.
Properties
Simple identities
Using the Kronecker delta one has,
and
if , or more generally if k > n.
Also
and
Similar relationships involving the Stirling numbers hold for the Bernoulli polynomials. Many relations for the Stirling numbers shadow similar relations on the binomial coefficients. The study of these 'shadow relationships' is termed umbral calculus and culminates in the theory of Sheffer sequences. Generalizations of the Stirling numbers of both kinds to arbitrary complex-valued inputs may be extended through the relations of these triangles to the Stirling convolution polynomials.
Note that all the combinatorial proofs above use either binomials or multinomials of .
Therefore if is prime, then:
for .
Expansions for fixed k
Since the Stirling numbers are the coefficients of a polynomial with roots 0, 1, ..., , one has by Vieta's formulas that
In other words, the Stirling numbers of the first kind are given by elementary symmetric polynomials evaluated at 0, 1, ..., . In this form, the simple identities given above take the form
and so on.
One may produce alternative forms for the Stirling numbers of the first kind with a similar approach preceded by some algebraic manipulation: since
it follows from Newton's formulas that one can expand the Stirling numbers of the first kind in terms of generalized harmonic numbers. This yields identities like
where Hn is the harmonic number and Hn(m) is the generalized harmonic number
These relations can be generalized to give
where w(n, m) is defined recursively in terms of the generalized harmonic numbers by
(Here δ is the Kronecker delta function and is the Pochhammer symbol.)
For fixed these weighted harmonic number expansions are generated by the generating function
where the notation means extraction of the coefficient of from the following formal power series (see the non-exponential Bell polynomials and section 3 of ).
More generally, sums related to these weighted harmonic number expansions of the Stirling numbers of the first kind can be defined through generalized zeta series transforms of generating functions.
One can also "invert" the relations for these Stirling numbers given in terms of the -order harmonic numbers to write the integer-order generalized harmonic numbers in terms of weighted sums of terms involving the Stirling numbers of the first kind. For example, when the second-order and third-order harmonic numbers are given by
More generally, one can invert the Bell polynomial generating function for the Stirling numbers expanded in terms of the -order harmonic numbers to obtain that for integers
Finite sums
Since permutations are partitioned by number of cycles, one has
The identities
and
can be proved by the techniques at
Stirling numbers and exponential generating functions#Stirling numbers of the first kind and Binomial coefficient#Ordinary generating functions.
The table in section 6.1 of Concrete Mathematics provides a plethora of generalized forms of finite sums involving the Stirling numbers. Several particular finite sums relevant to this article include
Additionally, if we define the second-order Eulerian numbers by the triangular recurrence relation
we arrive at the following identity related to the form of the Stirling convolution polynomials which can be employed to generalize both Stirling number triangles to arbitrary real, or complex-valued, values of the input :
Particular expansions of the previous identity lead to the following identities expanding the Stirling numbers of the first kind for the first few small values of :
Software tools for working with finite sums involving Stirling numbers and Eulerian numbers are provided by the RISC Stirling.m package utilities in Mathematica. Other software packages for guessing formulas for sequences (and polynomial sequence sums) involving Stirling numbers and other special triangles is available for both Mathematica and Sage here and here, respectively.
Congruences
The following congruence identity may be proved via a generating function-based approach:
More recent results providing Jacobi-type J-fractions that generate the single factorial function and generalized factorial-related products lead to other new congruence results for the Stirling numbers of the first kind.
For example, working modulo we can prove that
Where is the Iverson bracket.
and working modulo we can similarly prove that
More generally, for fixed integers if we define the ordered roots
then we may expand congruences for these Stirling numbers defined as the coefficients
in the following form where the functions, , denote fixed
polynomials of degree in for each , , and :
Section 6.2 of the reference cited above provides more explicit expansions related to these congruences for the -order harmonic numbers and for the generalized factorial products, .
Generating functions
A variety of identities may be derived by manipulating the generating function (see change of basis):
Using the equality
it follows that
and
This identity is valid for formal power series, and the sum converges in the complex plane for |z| < 1.
Other identities arise by exchanging the order of summation, taking derivatives, making substitutions for z or u, etc. For example, we may derive:
or
and
where and are the Riemann zeta function and the Hurwitz zeta function respectively, and even evaluate this integral
where is the gamma function. There also exist more complicated expressions for the zeta-functions involving the Stirling numbers. One, for example, has
This series generalizes Hasse's series for the Hurwitz zeta-function (we obtain Hasse's series by setting k=1).
Asymptotics
The next estimate given in terms of the Euler gamma constant applies:
For fixed we have the following estimate :
Explicit formula
It is well-known that we don't know any one-sum formula for Stirling numbers of the first kind. A two-sum formula can be obtained using one of the symmetric formulae for Stirling numbers in conjunction with the explicit formula for Stirling numbers of the second kind.
As discussed earlier, by Vieta's formulas, one getThe Stirling number s(n,n-p) can be found from the formula
where The sum is a sum over all partitions of p.
Another exact nested sum expansion for these Stirling numbers is computed by elementary symmetric polynomials corresponding to the coefficients in of a product of the form . In particular, we see that
Newton's identities combined with the above expansions may be used to give an alternate proof of the weighted expansions involving the generalized harmonic numbers already noted above.
Relations to natural logarithm function
The nth derivative of the μth power of the natural logarithm involves the signed Stirling numbers of the first kind:
where is the falling factorial, and is the signed Stirling number.
It can be proved by using mathematical induction.
Other formulas
Stirling numbers of the first kind appear in the formula for Gregory coefficients and in a finite sum identity involving Bell numbers
Infinite series involving the finite sums with the Stirling numbers often lead to the special functions. For example
and
or even
where γm are the Stieltjes constants and δm,0 represents the Kronecker delta function.
Notice that this last identity immediately implies relations between the polylogarithm functions, the Stirling number exponential generating functions given above, and the Stirling-number-based power series for the generalized Nielsen polylogarithm functions.
Generalizations
There are many notions of generalized Stirling numbers that may be defined (depending on application) in a number of differing combinatorial contexts. In so much as the Stirling numbers of the first kind correspond to the coefficients of the distinct polynomial expansions of the single factorial function, , we may extend this notion to define triangular recurrence relations for more general classes of products.
In particular, for any fixed arithmetic function and symbolic parameters , related generalized factorial products of the form
may be studied from the point of view of the classes of generalized Stirling numbers of the first kind defined by the following coefficients of the powers of in the expansions of and then by the next corresponding triangular recurrence relation:
These coefficients satisfy a number of analogous properties to those for the Stirling numbers of the first kind as well as recurrence relations and functional equations related to the f-harmonic numbers, .
One special case of these bracketed coefficients corresponding to allows us to expand the multiple factorial, or multifactorial functions as polynomials in .
The Stirling numbers of both kinds, the binomial coefficients, and the first and second-order Eulerian numbers are all defined by special cases of a triangular super-recurrence of the form
for integers and where whenever or . In this sense, the form of the Stirling numbers of the first kind may also be generalized by this parameterized super-recurrence for fixed scalars (not all zero).
| Mathematics | Combinatorics | null |
2229759 | https://en.wikipedia.org/wiki/Ensis | Ensis | Ensis is a genus of medium-sized edible saltwater clams, littoral bivalve molluscs in the family Pharidae. Ensis, or razor clams, are known in much of Scotland as spoots, for the spouts of water they eject while burrowing into the sand, when visible at low tide. This term may also colloquially include members of the genus Solen. Ensis magnus are known as bendies due to their slightly curved shell.
Description
The shells are long, narrow, and parallel-sided. This shape resembles a closed, old-fashioned straight razor (a cut-throat razor), or a closed jackknife (pocket knife) and sometimes these clams are known as razor shells or jackknives. The shells in these species are fragile and can easily be damaged when digging for these clams.
Ecology
Ensis species live in clean sand on exposed beaches. They are capable of digging very rapidly; see the description under the Atlantic jackknife clam. Some clammers catch jackknives by pouring salt on the characteristic keyhole-shaped breathing holes. The clam then tries to escape the salt by coming up out of its hole, and at this point it is possible to gently grab the shell and pull it out of the ground.
Species
Thirteen species are currently recognised:
Ensis californicus (Dall, 1899)
Ensis directus (Conrad, 1843) – Atlantic jackknife clam
Ensis ensis (Linnaeus, 1758)
Ensis goreensis (Clessin, 1888)
Ensis macha (Molina, 1782)
Ensis magnus (Schumacher, 1817)
Ensis megistus (Pilsbry & McGinty, 1943)
Ensis minor (Chenu, 1843)
Ensis myrae (Berry, 1954)
Ensis nitidus (Clessin, 1888)
Ensis siliqua (Linnaeus, 1758) – pod razor
Ensis terranovensis (Vierna & Martínez-Lage, 2012)
Ensis tropicalis (Hertlein & Strong, 1955)
| Biology and health sciences | Bivalvia | Animals |
2230558 | https://en.wikipedia.org/wiki/Diadectes | Diadectes | Diadectes (meaning crosswise-biter) is an extinct genus of large reptiliomorphs or synapsids that lived during the early Permian period (Artinskian-Kungurian stages of the Cisuralian epoch, between 290 and 272 million years ago). Diadectes was one of the first herbivorous tetrapods, and also one of the first fully terrestrial vertebrates to attain large size.
Description
Diadectes was a heavily built animal, up to long, with a thick-boned skull, heavy vertebrae and ribs, massive limb girdles, and short, robust limbs. The nature of the limbs and vertebrae clearly indicates a terrestrial animal. The rib cage was assumed to be barrel-shaped, but new fossils show the ribs were actually sticking out to the sides.
Paleobiology
It possesses some characteristics of reptilians and amphibians, combining a reptile-like skeleton with a more primitive, seymouriamorph-like skull. Diadectes has been classified as belonging to the sister group of the amniotes.
Among its primitive features, Diadectes has a large otic notch (a feature found in all labyrinthodonts, but not in reptiles) with an ossified tympanum. At the same time, its teeth show advanced specialisations for an herbivorous diet that are not found in any other type of early Permian animal. The eight front teeth are spatulate and peg-like, and served as incisors that were used to nip off mouthfuls of vegetation. The broad, blunt cheek teeth show extensive wear associated with occlusion, and would have functioned as molars, grinding up the food. It also had a partial secondary palate, which meant it could chew its food and breathe at the same time, something many even more advanced reptiles were unable to do.
These traits are likely adaptations related to the animals' high-fiber, herbivorous diet, and evolved independently of similar traits seen in some reptilian groups. Many of the reptile-like details of the postcranial skeleton are possibly related to carrying the substantial trunk; these may be independently derived traits on Diadectes and their relatives. Though very similar, they would be analogous rather than homologous to those of early amniotes such as pelycosaurs and pareiasaurs, as the first reptiles evolved from small, swamp-dwelling animals like Casineria and Westlothiana. The phenomenon of unrelated animals evolving similarly is known as convergent evolution.
Discovery
Diadectes was first named and described by American paleontologist Edward Drinker Cope in 1878, based on part of a lower jaw (AMNH 4360) from the Permian of Texas. Cope noted: "Teeth with short and much compressed crowns, whose long axis is transverse to that of the jaws," the feature expressed in the generic name Diadectes "crosswise biter" (from Greek dia "crosswise" + Greek dēktēs "biter"). He described the animal as "in all probability, herbivorous." Cope's neo-Latin type species name sideropelicus (from Greek sidēros "iron" + Greek pēlos "clay" + -ikos) "of iron clay" alluded to the Wichita beds in Texas, where the fossil was found.
Diadectes fossil remains are known from a number of locations across North America, especially the Texas Red Beds (Wichita and Clear Fork).
Classification and species
Numerous species have been assigned to Diadectes, though most of those have proven to be synonyms of one another. Similarly, many supposed separate genera of diadectids have been shown to be junior synonyms of Diadectes. One of these, Nothodon, was actually published by Othniel Charles Marsh five days before the name Diadectes was published by his rival Cope. Despite this fact, in 1912, Case synonymized the two names and treated Diadectes as the senior synonym, which has been followed by other paleontologists since, despite the fact that it violates the rules of International Code of Zoological Nomenclature (ICZN).
Phylogeny
A phylogenetic analysis of Diadectes and related diadectids was presented in an unpublished PhD thesis by Richard Kissel in 2010. Previous phylogenetic analyses of diadectids had found D. sanmiguelensis and D. absitus to be more basal than other species of Diadectes, outside the derived clade composed of these species. In these analyses, Diasparactus zenos was more closely related to the other species of Diadectes than was D. sanmiguelensis and D. absitus, making Diadectes paraphyletic. Kissel recovered this paraphyly in his analysis and proposed the new genus name "Oradectes" for D. sanmiguelensis, and "Silvadectes" for D. absitus. Below is the cladogram from Kissel's thesis:
However, according to the ICZN, a name presented in an initially unpublished thesis such as Kissel's is not valid. Because the names "Oradectes" and "Silvadectes" have not yet been formally erected in a published paper, they were not, as of 2010, considered valid. A 2024 paper formally erected the genus Kuwavaatakdectes for D. sanmiguelensis. The same paper also named another Diadectes species, D. dreigleichenensis, which coexisted with D. absitus in the Tambach Formation of Germany.
| Biology and health sciences | Reptiliomorphs | Animals |
2231059 | https://en.wikipedia.org/wiki/Superheavy%20element | Superheavy element | Superheavy elements, also known as transactinide elements, transactinides, or super-heavy elements, or superheavies for short, are the chemical elements with atomic number greater than 104. The superheavy elements are those beyond the actinides in the periodic table; the last actinide is lawrencium (atomic number 103). By definition, superheavy elements are also transuranium elements, i.e., having atomic numbers greater than that of uranium (92). Depending on the definition of group 3 adopted by authors, lawrencium may also be included to complete the 6d series.
Glenn T. Seaborg first proposed the actinide concept, which led to the acceptance of the actinide series. He also proposed a transactinide series ranging from element 104 to 121 and a superactinide series approximately spanning elements 122 to 153 (though more recent work suggests the end of the superactinide series to occur at element 157 instead). The transactinide seaborgium was named in his honor.
Superheavies are radioactive and have only been obtained synthetically in laboratories. No macroscopic sample of any of these elements has ever been produced. Superheavies are all named after physicists and chemists or important locations involved in the synthesis of the elements.
IUPAC defines an element to exist if its lifetime is longer than 10 seconds, which is the time it takes for the atom to form an electron cloud.
The known superheavies form part of the 6d and 7p series in the periodic table. Except for rutherfordium and dubnium (and lawrencium if it is included), even the longest-lived known isotopes of superheavies have half-lives of minutes or less. The element naming controversy involved elements 102–109. Some of these elements thus used systematic names for many years after their discovery was confirmed. (Usually the systematic names are replaced with permanent names proposed by the discoverers relatively soon after a discovery has been confirmed.)
Introduction
Synthesis of superheavy nuclei
A superheavy atomic nucleus is created in a nuclear reaction that combines two other nuclei of unequal size into one; roughly, the more unequal the two nuclei in terms of mass, the greater the possibility that the two react. The material made of the heavier nuclei is made into a target, which is then bombarded by the beam of lighter nuclei. Two nuclei can only fuse into one if they approach each other closely enough; normally, nuclei (all positively charged) repel each other due to electrostatic repulsion. The strong interaction can overcome this repulsion but only within a very short distance from a nucleus; beam nuclei are thus greatly accelerated in order to make such repulsion insignificant compared to the velocity of the beam nucleus. The energy applied to the beam nuclei to accelerate them can cause them to reach speeds as high as one-tenth of the speed of light. However, if too much energy is applied, the beam nucleus can fall apart.
Coming close enough alone is not enough for two nuclei to fuse: when two nuclei approach each other, they usually remain together for about 10 seconds and then part ways (not necessarily in the same composition as before the reaction) rather than form a single nucleus. This happens because during the attempted formation of a single nucleus, electrostatic repulsion tears apart the nucleus that is being formed. Each pair of a target and a beam is characterized by its cross section—the probability that fusion will occur if two nuclei approach one another expressed in terms of the transverse area that the incident particle must hit in order for the fusion to occur. This fusion may occur as a result of the quantum effect in which nuclei can tunnel through electrostatic repulsion. If the two nuclei can stay close past that phase, multiple nuclear interactions result in redistribution of energy and an energy equilibrium.
The resulting merger is an excited state—termed a compound nucleus—and thus it is very unstable. To reach a more stable state, the temporary merger may fission without formation of a more stable nucleus. Alternatively, the compound nucleus may eject a few neutrons, which would carry away the excitation energy; if the latter is not sufficient for a neutron expulsion, the merger would produce a gamma ray. This happens in about 10 seconds after the initial nuclear collision and results in creation of a more stable nucleus. The definition by the IUPAC/IUPAP Joint Working Party (JWP) states that a chemical element can only be recognized as discovered if a nucleus of it has not decayed within 10 seconds. This value was chosen as an estimate of how long it takes a nucleus to acquire electrons and thus display its chemical properties.
Decay and detection
The beam passes through the target and reaches the next chamber, the separator; if a new nucleus is produced, it is carried with this beam. In the separator, the newly produced nucleus is separated from other nuclides (that of the original beam and any other reaction products) and transferred to a surface-barrier detector, which stops the nucleus. The exact location of the upcoming impact on the detector is marked; also marked are its energy and the time of the arrival. The transfer takes about 10 seconds; in order to be detected, the nucleus must survive this long. The nucleus is recorded again once its decay is registered, and the location, the energy, and the time of the decay are measured.
Stability of a nucleus is provided by the strong interaction. However, its range is very short; as nuclei become larger, its influence on the outermost nucleons (protons and neutrons) weakens. At the same time, the nucleus is torn apart by electrostatic repulsion between protons, and its range is not limited. Total binding energy provided by the strong interaction increases linearly with the number of nucleons, whereas electrostatic repulsion increases with the square of the atomic number, i.e. the latter grows faster and becomes increasingly important for heavy and superheavy nuclei. Superheavy nuclei are thus theoretically predicted and have so far been observed to predominantly decay via decay modes that are caused by such repulsion: alpha decay and spontaneous fission. Almost all alpha emitters have over 210 nucleons, and the lightest nuclide primarily undergoing spontaneous fission has 238. In both decay modes, nuclei are inhibited from decaying by corresponding energy barriers for each mode, but they can be tunneled through.
Alpha particles are commonly produced in radioactive decays because the mass of an alpha particle per nucleon is small enough to leave some energy for the alpha particle to be used as kinetic energy to leave the nucleus. Spontaneous fission is caused by electrostatic repulsion tearing the nucleus apart and produces various nuclei in different instances of identical nuclei fissioning. As the atomic number increases, spontaneous fission rapidly becomes more important: spontaneous fission partial half-lives decrease by 23 orders of magnitude from uranium (element 92) to nobelium (element 102), and by 30 orders of magnitude from thorium (element 90) to fermium (element 100). The earlier liquid drop model thus suggested that spontaneous fission would occur nearly instantly due to disappearance of the fission barrier for nuclei with about 280 nucleons. The later nuclear shell model suggested that nuclei with about 300 nucleons would form an island of stability in which nuclei will be more resistant to spontaneous fission and will primarily undergo alpha decay with longer half-lives. Subsequent discoveries suggested that the predicted island might be further than originally anticipated; they also showed that nuclei intermediate between the long-lived actinides and the predicted island are deformed, and gain additional stability from shell effects. Experiments on lighter superheavy nuclei, as well as those closer to the expected island, have shown greater than previously anticipated stability against spontaneous fission, showing the importance of shell effects on nuclei.
Alpha decays are registered by the emitted alpha particles, and the decay products are easy to determine before the actual decay; if such a decay or a series of consecutive decays produces a known nucleus, the original product of a reaction can be easily determined. (That all decays within a decay chain were indeed related to each other is established by the location of these decays, which must be in the same place.) The known nucleus can be recognized by the specific characteristics of decay it undergoes such as decay energy (or more specifically, the kinetic energy of the emitted particle). Spontaneous fission, however, produces various nuclei as products, so the original nuclide cannot be determined from its daughters.
The information available to physicists aiming to synthesize a superheavy element is thus the information collected at the detectors: location, energy, and time of arrival of a particle to the detector, and those of its decay. The physicists analyze this data and seek to conclude that it was indeed caused by a new element and could not have been caused by a different nuclide than the one claimed. Often, provided data is insufficient for a conclusion that a new element was definitely created and there is no other explanation for the observed effects; errors in interpreting data have been made.
History
Early predictions
The heaviest element known at the end of the 19th century was uranium, with an atomic mass of about 240 (now known to be 238) amu. Accordingly, it was placed in the last row of the periodic table; this fueled speculation about the possible existence of elements heavier than uranium and why A = 240 seemed to be the limit. Following the discovery of the noble gases, beginning with argon in 1895, the possibility of heavier members of the group was considered. Danish chemist Julius Thomsen proposed in 1895 the existence of a sixth noble gas with Z = 86, A = 212 and a seventh with Z = 118, A = 292, the last closing a 32-element period containing thorium and uranium. In 1913, Swedish physicist Johannes Rydberg extended Thomsen's extrapolation of the periodic table to include even heavier elements with atomic numbers up to 460, but he did not believe that these superheavy elements existed or occurred in nature.
In 1914, German physicist Richard Swinne proposed that elements heavier than uranium, such as those around Z = 108, could be found in cosmic rays. He suggested that these elements may not necessarily have decreasing half-lives with increasing atomic number, leading to speculation about the possibility of some longer-lived elements at Z = 98–102 and Z = 108–110 (though separated by short-lived elements). Swinne published these predictions in 1926, believing that such elements might exist in Earth's core, iron meteorites, or the ice caps of Greenland where they had been locked up from their supposed cosmic origin.
Discoveries
Work performed from 1961 to 2013 at four labs – Lawrence Berkeley National Laboratory in the US, the Joint Institute for Nuclear Research in the USSR (later Russia), the GSI Helmholtz Centre for Heavy Ion Research in Germany, and Riken in Japan – identified and confirmed the elements lawrencium to oganesson according to the criteria of the IUPAC–IUPAP Transfermium Working Groups and subsequent Joint Working Parties. These discoveries complete the seventh row of the periodic table. The next two elements, ununennium (Z = 119) and unbinilium (Z = 120), have not yet been synthesized. They would begin an eighth period.
List of elements
103 Lawrencium, Lr, for Ernest Lawrence; sometimes but not always included
104 Rutherfordium, Rf, for Ernest Rutherford
105 Dubnium, Db, for the town of Dubna, near Moscow
106 Seaborgium, Sg, for Glenn T. Seaborg
107 Bohrium, Bh, for Niels Bohr
108 Hassium, Hs, for Hassia (Hesse), location of Darmstadt
109 Meitnerium, Mt, for Lise Meitner
110 Darmstadtium, Ds, for Darmstadt)
111 Roentgenium, Rg, for Wilhelm Röntgen
112 Copernicium, Cn, for Nicolaus Copernicus
113 Nihonium, Nh, for Nihon (Japan), location of the Riken institute
114 Flerovium, Fl, for Russian physicist Georgy Flyorov
115 Moscovium, Mc, for Moscow
116 Livermorium, Lv, for Lawrence Livermore National Laboratory
117 Tennessine, Ts, for Tennessee, location of Oak Ridge National Laboratory
118 Oganesson, Og, for Russian physicist Yuri Oganessian
Characteristics
Due to their short half-lives (for example, the most stable known isotope of seaborgium has a half-life of 14 minutes, and half-lives decrease with increasing atomic number) and the low yield of the nuclear reactions that produce them, new methods have had to be created to determine their gas-phase and solution chemistry based on very small samples of a few atoms each. Relativistic effects become very important in this region of the periodic table, causing the filled 7s orbitals, empty 7p orbitals, and filling 6d orbitals to all contract inward toward the atomic nucleus. This causes a relativistic stabilization of the 7s electrons and makes the 7p orbitals accessible in low excitation states.
Elements 103 to 112, lawrencium to copernicium, form the 6d series of transition elements. Experimental evidence shows that elements 103–108 behave as expected for their position in the periodic table, as heavier homologs of lutetium through osmium. They are expected to have ionic radii between those of their 5d transition metal homologs and their actinide pseudohomologs: for example, Rf is calculated to have ionic radius 76 pm, between the values for Hf (71 pm) and Th (94 pm). Their ions should also be less polarizable than those of their 5d homologs. Relativistic effects are expected to reach a maximum at the end of this series, at roentgenium (element 111) and copernicium (element 112). Nevertheless, many important properties of the transactinides are still not yet known experimentally, though theoretical calculations have been performed.
Elements 113 to 118, nihonium to oganesson, should form a 7p series, completing the seventh period in the periodic table. Their chemistry will be greatly influenced by the very strong relativistic stabilization of the 7s electrons and a strong spin–orbit coupling effect "tearing" the 7p subshell apart into two sections, one more stabilized (7p, holding two electrons) and one more destabilized (7p, holding four electrons). Lower oxidation states should be stabilized here, continuing group trends, as both the 7s and 7p electrons exhibit the inert-pair effect. These elements are expected to largely continue to follow group trends, though with relativistic effects playing an increasingly larger role. In particular, the large 7p splitting results in an effective shell closure at flerovium (element 114) and a hence much higher than expected chemical activity for oganesson (element 118).
Element 118 is the last element that has been synthesized. The next two elements, 119 and 120, should form an 8s series and be an alkali and alkaline earth metal respectively. The 8s electrons are expected to be relativistically stabilized, so that the trend toward higher reactivity down these groups will reverse and the elements will behave more like their period 5 homologs, rubidium and strontium. The 7p orbital is still relativistically destabilized, potentially giving these elements larger ionic radii and perhaps even being able to participate chemically. In this region, the 8p electrons are also relativistically stabilized, resulting in a ground-state 8s8p valence electron configuration for element 121. Large changes are expected to occur in the subshell structure in going from element 120 to element 121: for example, the radius of the 5g orbitals should drop drastically, from 25 Bohr units in element 120 in the excited [Og] 5g 8s configuration to 0.8 Bohr units in element 121 in the excited [Og] 5g 7d 8s configuration, in a phenomenon called "radial collapse". Element 122 should add either a further 7d or a further 8p electron to element 121's electron configuration. Elements 121 and 122 should be similar to actinium and thorium respectively.
At element 121, the superactinide series is expected to begin, when the 8s electrons and the filling 8p, 7d, 6f, and 5g subshells determine the chemistry of these elements. Complete and accurate calculations are not available for elements beyond 123 because of the extreme complexity of the situation: the 5g, 6f, and 7d orbitals should have about the same energy level, and in the region of element 160 the 9s, 8p, and 9p orbitals should also be about equal in energy. This will cause the electron shells to mix so that the block concept no longer applies very well, and will also result in novel chemical properties that will make positioning these elements in a periodic table very difficult.
Beyond superheavy elements
It has been suggested that elements beyond Z = 126 be called beyond superheavy elements. Other sources refer to elements around Z = 164 as hyperheavy elements.
| Physical sciences | Periods | Chemistry |
2232092 | https://en.wikipedia.org/wiki/Lutein | Lutein | Lutein (; from Latin luteus meaning "yellow") is a xanthophyll and one of 600 known naturally occurring carotenoids. Lutein is synthesized only by plants, and like other xanthophylls is found in high quantities in green leafy vegetables such as spinach, kale and yellow carrots. In green plants, xanthophylls act to modulate light energy and serve as non-photochemical quenching agents to deal with triplet chlorophyll, an excited form of chlorophyll which is overproduced at very high light levels during photosynthesis. See xanthophyll cycle for this topic.
Animals obtain lutein by ingesting plants. In the human retina, lutein is absorbed from blood specifically into the macula lutea, although its precise role in the body is unknown. Lutein is also found in egg yolks and animal fats.
Lutein is isomeric with zeaxanthin, differing only in the placement of one double bond. Lutein and zeaxanthin can be interconverted in the body through an intermediate called meso-zeaxanthin. The principal natural stereoisomer of lutein is (3R,3R,6R)-beta,epsilon-carotene-3,3-diol. Lutein is a lipophilic molecule and is generally insoluble in water. The presence of the long chromophore of conjugated double bonds (polyene chain) provides the distinctive light-absorbing properties. The polyene chain is susceptible to oxidative degradation by light or heat and is chemically unstable in acids.
Lutein is present in plants as fatty-acid esters, with one or two fatty acids bound to the two hydroxyl-groups. For this reason, saponification (de-esterification) of lutein esters to yield free lutein may yield lutein in any ratio from 1:1 to 1:2 molar ratio with the saponifying fatty acid.
As a pigment
This xanthophyll, like its sister compound zeaxanthin, has primarily been used in food and supplement manufacturing as a colorant due to its yellow-red color. Lutein absorbs blue light and therefore appears yellow at low concentrations and orange-red at high concentrations.
Many songbirds (like golden oriole, evening grosbeak, yellow warbler, common yellowthroat and Javan green magpies, but not American goldfinch or yellow canaries) deposit lutein obtained from the diet into growing tissues to color their feathers.
Role in human eyes
Although lutein is concentrated in the macula – a small area of the retina responsible for three-color vision – the precise functional role of retinal lutein has not been determined.
Macular degeneration
In 2013, findings of the Age-Related Eye Disease Study (AREDS2) showed that a dietary supplement formulation containing lutein reduced progression of age-related macular degeneration (AMD) by 25 percent. However, lutein and zeaxanthin had no overall effect on preventing AMD, but rather "the participants with low dietary intake of lutein and zeaxanthin at the start of the study, but who took an AREDS formulation with lutein and zeaxanthin during the study, were about 25 percent less likely to develop advanced AMD compared with participants with similar dietary intake who did not take lutein and zeaxanthin."
In AREDS2, participants took one of four AREDS formulations: the original AREDS formulation, AREDS formulation with no beta-carotene, AREDS with low zinc, AREDS with no beta-carotene and low zinc. In addition, they took one of four additional supplement or combinations including lutein and zeaxanthin (10 mg and 2 mg), omega-3 fatty acids (1,000 mg), lutein/zeaxanthin and omega-3 fatty acids, or placebo. The study reported that there was no overall additional benefit from adding omega-3 fatty acids or lutein and zeaxanthin to the formulation. However, the study did find benefits in two subgroups of participants: those not given beta-carotene, and those who had little lutein and zeaxanthin in their diets. Removing beta-carotene did not curb the formulation's protective effect against developing advanced AMD, which was important given that high doses of beta-carotene had been linked to higher risk of lung cancers in smokers. It was recommended to replace beta-carotene with lutein and zeaxanthin in future formulations for these reasons.
Three subsequent meta-analyses of dietary lutein and zeaxanthin concluded that these carotenoids lower the risk of progression from early stage AMD to late stage AMD.
An updated 2023 Cochrane review of 26 studies from several countries, however, concluded that dietary supplements containing zeaxanthin and lutein alone have little effect when compared to placebo on the progression of AMD. In general, there remains insufficient evidence to assess the effectiveness of dietary or supplemental zeaxanthin or lutein in treatment or prevention of early AMD.
Cataract research
There is preliminary epidemiological evidence that increasing lutein and zeaxanthin intake lowers the risk of cataract development. Consumption of more than 2.4 mg of lutein/zeaxanthin daily from foods and supplements was significantly correlated with reduced incidence of nuclear lens opacities, as revealed from data collected during a 13- to 15-year period in one study.
Two meta-analyses confirm a correlation between high diet content or high serum concentrations of lutein and zeaxanthin and a decrease in the risk of cataract. There is only one published clinical intervention trial testing for an effect of lutein and zeaxanthin supplementation on cataracts. The AREDS2 trial enrolled subjects at risk for progression to advanced age-related macular degeneration. Overall, the group getting lutein (10 mg) and zeaxanthin (2 mg) were NOT less likely to progress to needing cataract surgery. The authors speculated that there may be a cataract prevention benefit for people with low dietary intake of lutein and zeaxanthin, but recommended more research.
In diet
Lutein is a natural part of a human diet found in orange-yellow fruits and flowers, and in leafy vegetables. According to the NHANES 2013-2014 survey, adults in the United States consume on average 1.7 mg/day of lutein and zeaxanthin combined. No recommended dietary allowance currently exists for lutein. Some positive health effects have been seen at dietary intake levels of 6–10 mg/day. The only definitive side effect of excess lutein consumption is bronzing of the skin (carotenodermia).
As a food additive, lutein has the E number E161b (INS number 161b) and is extracted from the petals of African marigold (Tagetes erecta). It is approved for use in the EU and Australia and New Zealand. In the United States lutein may not be used as a food coloring for foods intended for human consumption, but can be added to animal feed and is allowed as a human dietary supplement often in combination with zeaxanthin. Example: lutein fed to chickens will show up in skin color and egg yolk color.
Some foods contain relatively high amounts of lutein:
Safety
In humans, the Observed Safe Level (OSL) for lutein, based on a non-government organization evaluation, is 20 mg/day. Although much higher levels have been tested without adverse effects and may also be safe, the data for intakes above the OSL are not sufficient for a confident conclusion of long-term safety. Neither the U.S. Food and Drug Administration nor the European Food Safety Authority considers lutein an essential nutrient or has acted to set a tolerable upper intake level.
Commercial value
The lutein market is segmented into pharmaceutical, dietary supplement, food, pet food, and animal and fish feed. The pharmaceutical market for lutein is estimated to be about US$190 million, and the nutraceutical and food categories are estimated to be about US$110 million. Pet food and other animal applications for lutein are estimated at US$175 million annually. This includes chickens (usually in combination with other carotenoids), to get color in egg yolks, and fish farms to color the flesh closer to wild-caught color. In the dietary supplement industry, the major market for lutein is for products with claims of helping maintain eye health. Newer applications are emerging in oral and topical products for skin health. Skin health via orally consumed supplements is one of the fastest growing areas of the US$2 billion carotenoid market.
| Biology and health sciences | Biological pigments | Biology |
5604894 | https://en.wikipedia.org/wiki/Presbyornis | Presbyornis | Presbyornis is an extinct genus of presbyornithid bird from North America during the Paleogene period, between the Late Paleocene and Early Eocene.
History of discovery
The fossil record of P. pervetus includes many complete skeletons from Green River Formation sites (Early Eocene), suggesting that the birds nested in colonies and that they possibly died due to volcanism or botulism, the latter of which is similar to many colony-nesting waterfowl or shorebirds today. Fossils identified as P. cf. pervetus have been discovered from the Margaret Formation of Ellesmere Island, where the remains of Gastornis sp. have also been found.
P. recurvirostra is known from a partial wing (KUVP 10105) found in the Colton Formation, from the Late Paleocene to Early Eocene sediments of the Wasatch Plateau near Ephraim, Utah. P. isoni, much larger than P. pervetus, is known from the Late Paleocene Aquia Formation in Maryland, based on the partial humerus (USNM 294116) and partial fingerbone (USNM 294117) that were initially described, as well as a complete humerus (SMM P96.9.2). Three humeri that were initially believed to be from Headonornis are suggested to belong to P. isoni, and the holotype coracoid of Headonornis may also be assigned to as P. isoni, though these claims require additional material for confirmation. However, Headonornis is now referred to as a stem group representative of the Phoenicopteriformes.
The holotype and paratypes of "P." mongoliensis are known from the Early Eocene of Mongolia, but these fragmentary specimens are poorly preserved and they likely belong to a stem Phoenicopterimorphae, not a presbyornithid. Undescribed fossils are also known from the Paleocene of Utah. Wunketru howardae, previously thought to be a species of Telmabates or a junior synonym of P. pervetus, is now considered a distinct anseriform.
Description
Along with Teviornis, Presbyornis was one of the earliest stem anseriforms. Because of its long legs and neck, Presbyornis could stand up to tall and was initially mistaken for a flamingo, but it was reclassified as an anseriform when the duck-like anatomy of its skull and bill was found. Later, it was believed to represent a transitional stage between the anseriforms and the shorebirds, but it is now considered a member of an extinct group of anseriforms which was most closely related to modern screamers. Judging from numerous fossil findings, Presbyornis is presumed to have lived in colonies around shallow lakes. Its broad, flat bill was used to filter food (small plants and animals) from the water, in the manner of today's dabbling ducks.
| Biology and health sciences | Prehistoric birds | Animals |
5605480 | https://en.wikipedia.org/wiki/Thermostability | Thermostability | In materials science and molecular biology, thermostability is the ability of a substance to resist irreversible change in its chemical or physical structure, often by resisting decomposition or polymerization, at a high relative temperature.
Thermostable materials may be used industrially as fire retardants. A thermostable plastic, an uncommon and unconventional term, is likely to refer to a thermosetting plastic that cannot be reshaped when heated, than to a thermoplastic that can be remelted and recast.
Thermostability is also a property of some proteins. To be a thermostable protein means to be resistant to changes in protein structure due to applied heat.
Thermostable proteins
Most life-forms on Earth live at temperatures of less than 50 °C, commonly from 15 to 50 °C. Within these organisms are macromolecules (proteins and nucleic acids) which form the three-dimensional structures essential to their enzymatic activity. Above the native temperature of the organism, thermal energy may cause the unfolding and denaturation, as the heat can disrupt the intramolecular bonds in the tertiary and quaternary structure. This unfolding will result in loss in enzymatic activity, which is understandably deleterious to continuing life-functions. An example of such is the denaturing of proteins in albumen from a clear, nearly colourless liquid to an opaque white, insoluble gel.
Proteins capable of withstanding such high temperatures compared to proteins that cannot, are generally from microorganisms that are hyperthermophiles. Such organisms can withstand above 50 °C temperatures as they usually live within environments of 85 °C and above. Certain thermophilic life-forms exist which can withstand temperatures above this, and have corresponding adaptations to preserve protein function at these temperatures. These can include altered bulk properties of the cell to stabilize all proteins, and specific changes to individual proteins. Comparing homologous proteins present in these thermophiles and other organisms reveal some differences in the protein structure. One notable difference is the presence of extra hydrogen bonds in the thermophile's proteins—meaning that the protein structure is more resistant to unfolding. Similarly, thermostable proteins are rich in salt bridges or/and extra disulfide bridges stabilizing the structure. Other factors of protein thermostability are compactness of protein structure, oligomerization, and strength interaction between subunits.
Uses and applications
Polymerase chain reactions
Thermostable DNA polymerases such as Taq polymerase and Pfu DNA polymerase are used in polymerase chain reactions (PCR) where temperatures of 94 °C or over are used to melt DNA strands in the denaturation step of PCR. This resistance to high temperature allows for DNA polymerase to elongate DNA with a desired sequence of interest with the presence of dNTPs.
Feed additives
Enzymes are often added to animal feed to improve the health and growth of farmed animals, particularly chickens and pigs. The feed is typically treated with high pressure steam to kill bacteria such as Salmonella. Therefore the added enzymes (e.g. phytase and xylanase) must be able to withstand this thermal challenge without being irreversibly inactivated.
Protein purification
Knowledge of an enzyme's resistance to high temperatures is especially beneficial in protein purification. In the procedure of heat denaturation, one can subject a mixture of proteins to high temperatures, which will result in the denaturation of proteins that are not thermostable, and the isolation of the protein that is thermodynamically stable. One notable example of this is found in the purification of alkaline phosphatase from the hyperthermophile Pyrococcus abyssi. This enzyme is known for being heat stable at temperatures greater than 95 °C, and therefore can be partially purified by heating when heterologously expressed in E. coli. The increase in temperature causes the E. coli proteins to precipitate, while the P. abyssi alkaline phosphatase remains stably in solution.
Glycoside hydrolases
Another important group of thermostable enzymes are glycoside hydrolases. These enzymes are responsible of the degradation of the major fraction of biomass, the polysaccharides present in starch and lignocellulose. Thus, glycoside hydrolases are gaining great interest in biorefining applications in the future bioeconomy. Some examples are the production of monosaccharides for food applications as well as use as carbon source for microbial conversion in fuels (ethanol) and chemical intermediates, production of oligosaccharides for prebiotic applications and production of surfactants alkyl glycoside type. All of these processes often involve thermal treatments to facilitate the polysaccharide hydrolysis, hence give thermostable variants of glycoside hydrolases an important role in this context.
Approaches to improve thermostability of proteins
Protein engineering can be used to enhance the thermostability of proteins. A number of site-directed and random mutagenesis techniques, in addition to directed evolution, have been used to increase the thermostability of target proteins. Comparative methods have been used to increase the stability of mesophilic proteins based on comparison to thermophilic homologs. Additionally, analysis of the protein unfolding by molecular dynamics can be used to understand the process of unfolding and then design stabilizing mutations. Rational protein engineering for increasing protein thermostability includes mutations which truncate loops, increase salt bridges or hydrogen bonds, introduced disulfide bonds. In addition, ligand binding can increase the stability of the protein, particularly when purified. There are various different forces that allow for the thermostability of a particular protein. These forces include hydrophobic interactions, electrostatic interactions, and the presence of disulfide bonds. The overall amount of hydrophobicity present in a particular protein is responsible for its thermostability. Another type of force that is responsible for thermostability of a protein is the electrostatic interactions between molecules. These interactions include salt bridges and hydrogen bonds. Salt bridges are unaffected by high temperatures, therefore, are necessary for protein and enzyme stability. A third force used to increase thermostability in proteins and enzymes is the presence of disulfide bonds. They present covalent cross-linkages between the polypeptide chains. These bonds are the strongest because they're covalent bonds, making them stronger than intermolecular forces. Glycosylation is another way to improve the thermostability of proteins. Stereoelectronic effects in stabilizing interactions between carbohydrate and protein can lead to the thermostabilization of the glycosylated protein.
Cyclizing enzymes by covalently linking the N-terminus to the C-terminus has been applied to increase the thermostability of many enzymes. Intein cyclization and SpyTag/SpyCatcher cyclization have often been employed.
Thermostable toxins
Certain poisonous fungi contain thermostable toxins, such as amatoxin found in the death cap and autumn skullcap mushrooms and patulin from molds. Therefore, applying heat to these will not remove the toxicity and is of particular concern for food safety.
| Biology and health sciences | Basics | Biology |
14640471 | https://en.wikipedia.org/wiki/Mars | Mars | Mars is the fourth planet from the Sun. The surface of Mars is orange-red because it is covered in iron(III) oxide dust, giving it the nickname "the Red Planet". Mars is among the brightest objects in Earth's sky, and its high-contrast albedo features have made it a common subject for telescope viewing. It is classified as a terrestrial planet and is the second smallest of the Solar System's planets with a diameter of . In terms of orbital motion, a Martian solar day (sol) is equal to 24.6 hours, and a Martian solar year is equal to 1.88 Earth years (687 Earth days). Mars has two natural satellites that are small and irregular in shape: Phobos and Deimos.
The relatively flat plains in northern parts of Mars strongly contrast with the cratered terrain in southern highlands – this terrain observation is known as the Martian dichotomy. Mars hosts many enormous extinct volcanoes (the tallest is Olympus Mons, tall) and one of the largest canyons in the Solar System (Valles Marineris, long). Geologically, the planet is fairly active with marsquakes trembling underneath the ground, dust devils sweeping across the landscape, and cirrus clouds. Carbon dioxide is substantially present in Mars's polar ice caps and thin atmosphere. During a year, there are large surface temperature swings on the surface between to similar to Earth's seasons, as both planets have significant axial tilt.
Mars was formed approximately 4.5 billion years ago. During the Noachian period (4.5 to 3.5 billion years ago), Mars's surface was marked by meteor impacts, valley formation, erosion, and the possible presence of water oceans. The Hesperian period (3.5 to 3.3–2.9 billion years ago) was dominated by widespread volcanic activity and flooding that carved immense outflow channels. The Amazonian period, which continues to the present, has been marked by the wind as a dominant influence on geological processes. Due to Mars's geological history, the possibility of past or present life on Mars remains of great scientific interest.
Since the late 20th century, Mars has been explored by uncrewed spacecraft and rovers, with the first flyby by the Mariner 4 probe in 1965, the first orbit by the Mars 2 probe in 1971, and the first landing by the Viking 1 probe in 1976. As of 2023, there are at least 11 active probes orbiting Mars or on the Martian surface. Mars is an attractive target for future human exploration missions, though in the 2020s no such mission is planned.
Natural history
Scientists have theorized that during the Solar System's formation, Mars was created as the result of a random process of run-away accretion of material from the protoplanetary disk that orbited the Sun. Mars has many distinctive chemical features caused by its position in the Solar System. Elements with comparatively low boiling points, such as chlorine, phosphorus, and sulfur, are much more common on Mars than on Earth; these elements were probably pushed outward by the young Sun's energetic solar wind.
After the formation of the planets, the inner Solar System may have been subjected to the so-called Late Heavy Bombardment. About 60% of the surface of Mars shows a record of impacts from that era, whereas much of the remaining surface is probably underlain by immense impact basins caused by those events. However, more recent modeling has disputed the existence of the Late Heavy Bombardment. There is evidence of an enormous impact basin in the Northern Hemisphere of Mars, spanning , or roughly four times the size of the Moon's South Pole–Aitken basin, which would be the largest impact basin yet discovered if confirmed. It has been hypothesized that the basin was formed when Mars was struck by a Pluto-sized body about four billion years ago. The event, thought to be the cause of the Martian hemispheric dichotomy, created the smooth Borealis basin that covers 40% of the planet.
A 2023 study shows evidence, based on the orbital inclination of Deimos (a small moon of Mars), that Mars may once have had a ring system 3.5 billion years to 4 billion years ago. This ring system may have been formed from a moon, 20 times more massive than Phobos, orbiting Mars billions of years ago; and Phobos would be a remnant of that ring.
The geological history of Mars can be split into many periods, but the following are the three primary periods:
Noachian period: Formation of the oldest extant surfaces of Mars, 4.5 to 3.5 billion years ago. Noachian age surfaces are scarred by many large impact craters. The Tharsis bulge, a volcanic upland, is thought to have formed during this period, with extensive flooding by liquid water late in the period. Named after Noachis Terra.
Hesperian period: 3.5 to between 3.3 and 2.9 billion years ago. The Hesperian period is marked by the formation of extensive lava plains. Named after Hesperia Planum.
Amazonian period: between 3.3 and 2.9 billion years ago to the present. Amazonian regions have few meteorite impact craters but are otherwise quite varied. Olympus Mons formed during this period, with lava flows elsewhere on Mars. Named after Amazonis Planitia.
Geological activity is still taking place on Mars. The Athabasca Valles is home to sheet-like lava flows created about 200 million years ago. Water flows in the grabens called the Cerberus Fossae occurred less than 20 million years ago, indicating equally recent volcanic intrusions. The Mars Reconnaissance Orbiter has captured images of avalanches.
Physical characteristics
Mars is approximately half the diameter of Earth, with a surface area only slightly less than the total area of Earth's dry land. Mars is less dense than Earth, having about 15% of Earth's volume and 11% of Earth's mass, resulting in about 38% of Earth's surface gravity. Mars is the only presently known example of a desert planet, a rocky planet with a surface akin to that of Earth's hot deserts. The red-orange appearance of the Martian surface is caused by rust. It can look like butterscotch; other common surface colors include golden, brown, tan, and greenish, depending on the minerals present.
Internal structure
Like Earth, Mars is differentiated into a dense metallic core overlaid by less dense rocky layers. The outermost layer is the crust, which is on average about thick, with a minimum thickness of in Isidis Planitia, and a maximum thickness of in the southern Tharsis plateau. For comparison, Earth's crust averages 27.3 ± 4.8 km in thickness. The most abundant elements in the Martian crust are silicon, oxygen, iron, magnesium, aluminium, calcium, and potassium. Mars is confirmed to be seismically active; in 2019 it was reported that InSight had detected and recorded over 450 marsquakes and related events.
Beneath the crust is a silicate mantle responsible for many of the tectonic and volcanic features on the planet's surface. The upper Martian mantle is a low-velocity zone, where the velocity of seismic waves is lower than surrounding depth intervals. The mantle appears to be rigid down to the depth of about 250 km, giving Mars a very thick lithosphere compared to Earth. Below this the mantle gradually becomes more ductile, and the seismic wave velocity starts to grow again. The Martian mantle does not appear to have a thermally insulating layer analogous to Earth's lower mantle; instead, below 1050 km in depth, it becomes mineralogically similar to Earth's transition zone. At the bottom of the mantle lies a basal liquid silicate layer approximately 150–180 km thick.
Mars's iron and nickel core is completely molten, with no solid inner core. It is around half of Mars's radius, approximately 1650–1675 km, and is enriched in light elements such as sulfur, oxygen, carbon, and hydrogen.
Surface geology
Mars is a terrestrial planet with a surface that consists of minerals containing silicon and oxygen, metals, and other elements that typically make up rock. The Martian surface is primarily composed of tholeiitic basalt, although parts are more silica-rich than typical basalt and may be similar to andesitic rocks on Earth, or silica glass. Regions of low albedo suggest concentrations of plagioclase feldspar, with northern low albedo regions displaying higher than normal concentrations of sheet silicates and high-silicon glass. Parts of the southern highlands include detectable amounts of high-calcium pyroxenes. Localized concentrations of hematite and olivine have been found. Much of the surface is deeply covered by finely grained iron(III) oxide dust.
Although Mars has no evidence of a structured global magnetic field, observations show that parts of the planet's crust have been magnetized, suggesting that alternating polarity reversals of its dipole field have occurred in the past. This paleomagnetism of magnetically susceptible minerals is similar to the alternating bands found on Earth's ocean floors. One hypothesis, published in 1999 and re-examined in October 2005 (with the help of the Mars Global Surveyor), is that these bands suggest plate tectonic activity on Mars four billion years ago, before the planetary dynamo ceased to function and the planet's magnetic field faded.
The Phoenix lander returned data showing Martian soil to be slightly alkaline and containing elements such as magnesium, sodium, potassium and chlorine. These nutrients are found in soils on Earth. They are necessary for growth of plants. Experiments performed by the lander showed that the Martian soil has a basic pH of 7.7, and contains 0.6% perchlorate by weight, concentrations that are toxic to humans.
Streaks are common across Mars and new ones appear frequently on steep slopes of craters, troughs, and valleys. The streaks are dark at first and get lighter with age. The streaks can start in a tiny area, then spread out for hundreds of metres. They have been seen to follow the edges of boulders and other obstacles in their path. The commonly accepted hypotheses include that they are dark underlying layers of soil revealed after avalanches of bright dust or dust devils. Several other explanations have been put forward, including those that involve water or even the growth of organisms.
Environmental radiation levels on the surface are on average 0.64 millisieverts of radiation per day, and significantly less than the radiation of 1.84 millisieverts per day or 22 millirads per day during the flight to and from Mars. For comparison the radiation levels in low Earth orbit, where Earth's space stations orbit, are around 0.5 millisieverts of radiation per day. Hellas Planitia has the lowest surface radiation at about 0.342 millisieverts per day, featuring lava tubes southwest of Hadriacus Mons with potentially levels as low as 0.064 millisieverts per day, comparable to radiation levels during flights on Earth.
Geography and features
Although better remembered for mapping the Moon, Johann Heinrich von Mädler and Wilhelm Beer were the first areographers. They began by establishing that most of Mars's surface features were permanent and by more precisely determining the planet's rotation period. In 1840, Mädler combined ten years of observations and drew the first map of Mars.
Features on Mars are named from a variety of sources. Albedo features are named for classical mythology. Craters larger than roughly 50 km are named for deceased scientists and writers and others who have contributed to the study of Mars. Smaller craters are named for towns and villages of the world with populations of less than 100,000. Large valleys are named for the word "Mars" or "star" in various languages; smaller valleys are named for rivers.
Large albedo features retain many of the older names but are often updated to reflect new knowledge of the nature of the features. For example, Nix Olympica (the snows of Olympus) has become Olympus Mons (Mount Olympus). The surface of Mars as seen from Earth is divided into two kinds of areas, with differing albedo. The paler plains covered with dust and sand rich in reddish iron oxides were once thought of as Martian "continents" and given names like Arabia Terra (land of Arabia) or Amazonis Planitia (Amazonian plain). The dark features were thought to be seas, hence their names Mare Erythraeum, Mare Sirenum and Aurorae Sinus. The largest dark feature seen from Earth is Syrtis Major Planum. The permanent northern polar ice cap is named Planum Boreum. The southern cap is called Planum Australe.
Mars's equator is defined by its rotation, but the location of its Prime Meridian was specified, as was Earth's (at Greenwich), by choice of an arbitrary point; Mädler and Beer selected a line for their first maps of Mars in 1830. After the spacecraft Mariner 9 provided extensive imagery of Mars in 1972, a small crater (later called Airy-0), located in the Sinus Meridiani ("Middle Bay" or "Meridian Bay"), was chosen by Merton E. Davies, Harold Masursky, and Gérard de Vaucouleurs for the definition of 0.0° longitude to coincide with the original selection.
Because Mars has no oceans, and hence no "sea level", a zero-elevation surface had to be selected as a reference level; this is called the areoid of Mars, analogous to the terrestrial geoid. Zero altitude was defined by the height at which there is of atmospheric pressure. This pressure corresponds to the triple point of water, and it is about 0.6% of the sea level surface pressure on Earth (0.006 atm).
For mapping purposes, the United States Geological Survey divides the surface of Mars into thirty cartographic quadrangles, each named for a classical albedo feature it contains. In April 2023, The New York Times reported an updated global map of Mars based on images from the Hope spacecraft. A related, but much more detailed, global Mars map was released by NASA on 16 April 2023.
Volcanoes
The vast upland region Tharsis contains several massive volcanoes, which include the shield volcano Olympus Mons. The edifice is over wide. Because the mountain is so large, with complex structure at its edges, giving a definite height to it is difficult. Its local relief, from the foot of the cliffs which form its northwest margin to its peak, is over , a little over twice the height of Mauna Kea as measured from its base on the ocean floor. The total elevation change from the plains of Amazonis Planitia, over to the northwest, to the summit approaches , roughly three times the height of Mount Everest, which in comparison stands at just over . Consequently, Olympus Mons is either the tallest or second-tallest mountain in the Solar System; the only known mountain which might be taller is the Rheasilvia peak on the asteroid Vesta, at .
Impact topography
The dichotomy of Martian topography is striking: northern plains flattened by lava flows contrast with the southern highlands, pitted and cratered by ancient impacts. It is possible that, four billion years ago, the Northern Hemisphere of Mars was struck by an object one-tenth to two-thirds the size of Earth's Moon. If this is the case, the Northern Hemisphere of Mars would be the site of an impact crater in size, or roughly the area of Europe, Asia, and Australia combined, surpassing Utopia Planitia and the Moon's South Pole–Aitken basin as the largest impact crater in the Solar System.
Mars is scarred by a number of impact craters: a total of 43,000 observed craters with a diameter of or greater have been found. The largest exposed crater is Hellas, which is wide and deep, and is a light albedo feature clearly visible from Earth. There are other notable impact features, such as Argyre, which is around in diameter, and Isidis, which is around in diameter. Due to the smaller mass and size of Mars, the probability of an object colliding with the planet is about half that of Earth. Mars is located closer to the asteroid belt, so it has an increased chance of being struck by materials from that source. Mars is more likely to be struck by short-period comets, i.e., those that lie within the orbit of Jupiter.
Martian craters can have a morphology that suggests the ground became wet after the meteor impact.
Tectonic sites
The large canyon, Valles Marineris (Latin for 'Mariner Valleys, also known as Agathodaemon in the old canal maps), has a length of and a depth of up to . The length of Valles Marineris is equivalent to the length of Europe and extends across one-fifth the circumference of Mars. By comparison, the Grand Canyon on Earth is only long and nearly deep. Valles Marineris was formed due to the swelling of the Tharsis area, which caused the crust in the area of Valles Marineris to collapse. In 2012, it was proposed that Valles Marineris is not just a graben, but a plate boundary where of transverse motion has occurred, making Mars a planet with possibly a two-tectonic plate arrangement.
Holes and caves
Images from the Thermal Emission Imaging System (THEMIS) aboard NASA's Mars Odyssey orbiter have revealed seven possible cave entrances on the flanks of the volcano Arsia Mons. The caves, named after loved ones of their discoverers, are collectively known as the "seven sisters". Cave entrances measure from wide and they are estimated to be at least deep. Because light does not reach the floor of most of the caves, they may extend much deeper than these lower estimates and widen below the surface. "Dena" is the only exception; its floor is visible and was measured to be deep. The interiors of these caverns may be protected from micrometeoroids, UV radiation, solar flares and high energy particles that bombard the planet's surface.
Atmosphere
Mars lost its magnetosphere 4 billion years ago, possibly because of numerous asteroid strikes, so the solar wind interacts directly with the Martian ionosphere, lowering the atmospheric density by stripping away atoms from the outer layer. Both Mars Global Surveyor and Mars Express have detected ionized atmospheric particles trailing off into space behind Mars, and this atmospheric loss is being studied by the MAVEN orbiter. Compared to Earth, the atmosphere of Mars is quite rarefied. Atmospheric pressure on the surface today ranges from a low of on Olympus Mons to over in Hellas Planitia, with a mean pressure at the surface level of . The highest atmospheric density on Mars is equal to that found above Earth's surface. The resulting mean surface pressure is only 0.6% of Earth's . The scale height of the atmosphere is about , which is higher than Earth's , because the surface gravity of Mars is only about 38% of Earth's.
The atmosphere of Mars consists of about 96% carbon dioxide, 1.93% argon and 1.89% nitrogen along with traces of oxygen and water. The atmosphere is quite dusty, containing particulates about 1.5 μm in diameter which give the Martian sky a tawny color when seen from the surface. It may take on a pink hue due to iron oxide particles suspended in it. The concentration of methane in the Martian atmosphere fluctuates from about 0.24 ppb during the northern winter to about 0.65 ppb during the summer. Estimates of its lifetime range from 0.6 to 4 years, so its presence indicates that an active source of the gas must be present. Methane could be produced by non-biological process such as serpentinization involving water, carbon dioxide, and the mineral olivine, which is known to be common on Mars, or by Martian life.
Compared to Earth, its higher concentration of atmospheric CO2 and lower surface pressure may be why sound is attenuated more on Mars, where natural sources are rare apart from the wind. Using acoustic recordings collected by the Perseverance rover, researchers concluded that the speed of sound there is approximately 240 m/s for frequencies below 240 Hz, and 250 m/s for those above.
Auroras have been detected on Mars. Because Mars lacks a global magnetic field, the types and distribution of auroras there differ from those on Earth; rather than being mostly restricted to polar regions as is the case on Earth, a Martian aurora can encompass the planet. In September 2017, NASA reported radiation levels on the surface of the planet Mars were temporarily doubled, and were associated with an aurora 25 times brighter than any observed earlier, due to a massive, and unexpected, solar storm in the middle of the month.
Climate
Mars has seasons, alternating between its northern and southern hemispheres, similar to on Earth. Additionally the orbit of Mars has, compared to Earth's, a large eccentricity and approaches perihelion when it is summer in its southern hemisphere and winter in its northern, and aphelion when it is winter in its southern hemisphere and summer in its northern. As a result, the seasons in its southern hemisphere are more extreme and the seasons in its northern are milder than would otherwise be the case. The summer temperatures in the south can be warmer than the equivalent summer temperatures in the north by up to .
Martian surface temperatures vary from lows of about to highs of up to in equatorial summer. The wide range in temperatures is due to the thin atmosphere which cannot store much solar heat, the low atmospheric pressure (about 1% that of the atmosphere of Earth), and the low thermal inertia of Martian soil. The planet is 1.52 times as far from the Sun as Earth, resulting in just 43% of the amount of sunlight.
Mars has the largest dust storms in the Solar System, reaching speeds of over . These can vary from a storm over a small area, to gigantic storms that cover the entire planet. They tend to occur when Mars is closest to the Sun, and have been shown to increase global temperature.
Seasons also produce dry ice covering polar ice caps. Large areas of the polar regions of Mars
Hydrology
While Mars contains water in larger amounts, most of it is dust covered water ice at the Martian polar ice caps.
The volume of water ice in the south polar ice cap, if melted, would be enough to cover most of the surface of the planet with a depth of .
Water in its liquid form cannot prevail on the surface of Mars due to the low atmospheric pressure on Mars, which is less than 1% that of Earth, only at the lowest of elevations pressure and temperature is high enough for water being able to be liquid for short periods.
Water in the atmosphere is small, but enough to produce larger clouds of water ice and different cases of snow and frost, often mixed with snow of carbon dioxide dry ice.
Past hydrosphere
Landforms visible on Mars strongly suggest that liquid water has existed on the planet's surface. Huge linear swathes of scoured ground, known as outflow channels, cut across the surface in about 25 places. These are thought to be a record of erosion caused by the catastrophic release of water from subsurface aquifers, though some of these structures have been hypothesized to result from the action of glaciers or lava. One of the larger examples, Ma'adim Vallis, is long, much greater than the Grand Canyon, with a width of and a depth of in places. It is thought to have been carved by flowing water early in Mars's history. The youngest of these channels is thought to have formed only a few million years ago.
Elsewhere, particularly on the oldest areas of the Martian surface, finer-scale, dendritic networks of valleys are spread across significant proportions of the landscape. Features of these valleys and their distribution strongly imply that they were carved by runoff resulting from precipitation in early Mars history. Subsurface water flow and groundwater sapping may play important subsidiary roles in some networks, but precipitation was probably the root cause of the incision in almost all cases.
Along craters and canyon walls, there are thousands of features that appear similar to terrestrial gullies. The gullies tend to be in the highlands of the Southern Hemisphere and face the Equator; all are poleward of 30° latitude. A number of authors have suggested that their formation process involves liquid water, probably from melting ice, although others have argued for formation mechanisms involving carbon dioxide frost or the movement of dry dust. No partially degraded gullies have formed by weathering and no superimposed impact craters have been observed, indicating that these are young features, possibly still active. Other geological features, such as deltas and alluvial fans preserved in craters, are further evidence for warmer, wetter conditions at an interval or intervals in earlier Mars history. Such conditions necessarily require the widespread presence of crater lakes across a large proportion of the surface, for which there is independent mineralogical, sedimentological and geomorphological evidence. Further evidence that liquid water once existed on the surface of Mars comes from the detection of specific minerals such as hematite and goethite, both of which sometimes form in the presence of water.
History of observations and findings of water evidence
In 2004, Opportunity detected the mineral jarosite. This forms only in the presence of acidic water, showing that water once existed on Mars. The Spirit rover found concentrated deposits of silica in 2007 that indicated wet conditions in the past, and in December 2011, the mineral gypsum, which also forms in the presence of water, was found on the surface by NASA's Mars rover Opportunity. It is estimated that the amount of water in the upper mantle of Mars, represented by hydroxyl ions contained within Martian minerals, is equal to or greater than that of Earth at 50–300 parts per million of water, which is enough to cover the entire planet to a depth of .
On 18 March 2013, NASA reported evidence from instruments on the Curiosity rover of mineral hydration, likely hydrated calcium sulfate, in several rock samples including the broken fragments of "Tintina" rock and "Sutton Inlier" rock as well as in veins and nodules in other rocks like "Knorr" rock and "Wernicke" rock. Analysis using the rover's DAN instrument provided evidence of subsurface water, amounting to as much as 4% water content, down to a depth of , during the rover's traverse from the Bradbury Landing site to the Yellowknife Bay area in the Glenelg terrain. In September 2015, NASA announced that they had found strong evidence of hydrated brine flows in recurring slope lineae, based on spectrometer readings of the darkened areas of slopes. These streaks flow downhill in Martian summer, when the temperature is above −23 °C, and freeze at lower temperatures. These observations supported earlier hypotheses, based on timing of formation and their rate of growth, that these dark streaks resulted from water flowing just below the surface. However, later work suggested that the lineae may be dry, granular flows instead, with at most a limited role for water in initiating the process. A definitive conclusion about the presence, extent, and role of liquid water on the Martian surface remains elusive.
Researchers suspect much of the low northern plains of the planet were covered with an ocean hundreds of meters deep, though this theory remains controversial. In March 2015, scientists stated that such an ocean might have been the size of Earth's Arctic Ocean. This finding was derived from the ratio of protium to deuterium in the modern Martian atmosphere compared to that ratio on Earth. The amount of Martian deuterium (D/H = 9.3 ± 1.7 10−4) is five to seven times the amount on Earth (D/H = 1.56 10−4), suggesting that ancient Mars had significantly higher levels of water. Results from the Curiosity rover had previously found a high ratio of deuterium in Gale Crater, though not significantly high enough to suggest the former presence of an ocean. Other scientists caution that these results have not been confirmed, and point out that Martian climate models have not yet shown that the planet was warm enough in the past to support bodies of liquid water. Near the northern polar cap is the wide Korolev Crater, which the Mars Express orbiter found to be filled with approximately of water ice.
In November 2016, NASA reported finding a large amount of underground ice in the Utopia Planitia region. The volume of water detected has been estimated to be equivalent to the volume of water in Lake Superior (which is 12,100 cubic kilometers). During observations from 2018 through 2021, the ExoMars Trace Gas Orbiter spotted indications of water, probably subsurface ice, in the Valles Marineris canyon system.
Orbital motion
Mars's average distance from the Sun is roughly , and its orbital period is 687 (Earth) days. The solar day (or sol) on Mars is only slightly longer than an Earth day: 24 hours, 39 minutes, and 35.244 seconds. A Martian year is equal to 1.8809 Earth years, or 1 year, 320 days, and 18.2 hours. The gravitational potential difference and thus the delta-v needed to transfer between Mars and Earth is the second lowest for Earth.
The axial tilt of Mars is 25.19° relative to its orbital plane, which is similar to the axial tilt of Earth. As a result, Mars has seasons like Earth, though on Mars they are nearly twice as long because its orbital period is that much longer. In the present day, the orientation of the north pole of Mars is close to the star Deneb.
Mars has a relatively pronounced orbital eccentricity of about 0.09; of the seven other planets in the Solar System, only Mercury has a larger orbital eccentricity. It is known that in the past, Mars has had a much more circular orbit. At one point, 1.35 million Earth years ago, Mars had an eccentricity of roughly 0.002, much less than that of Earth today. Mars's cycle of eccentricity is 96,000 Earth years compared to Earth's cycle of 100,000 years.
Mars has its closest approach to Earth (opposition) in a synodic period of 779.94 days. It should not be confused with Mars conjunction, where the Earth and Mars are at opposite sides of the Solar System and form a straight line crossing the Sun. The average time between the successive oppositions of Mars, its synodic period, is 780 days; but the number of days between successive oppositions can range from 764 to 812. The distance at close approach varies between about due to the planets' elliptical orbits, which causes comparable variation in angular size. At their furthest Mars and Earth can be as far as apart. Mars comes into opposition from Earth every 2.1 years. The planets come into opposition near Mars's perihelion in 2003, 2018 and 2035, with the 2020 and 2033 events being particularly close to perihelic opposition.
The mean apparent magnitude of Mars is +0.71 with a standard deviation of 1.05. Because the orbit of Mars is eccentric, the magnitude at opposition from the Sun can range from about −3.0 to −1.4. The minimum brightness is magnitude +1.86 when the planet is near aphelion and in conjunction with the Sun. At its brightest, Mars (along with Jupiter) is second only to Venus in apparent brightness. Mars usually appears distinctly yellow, orange, or red. When farthest away from Earth, it is more than seven times farther away than when it is closest. Mars is usually close enough for particularly good viewing once or twice at 15-year or 17-year intervals. Optical ground-based telescopes are typically limited to resolving features about across when Earth and Mars are closest because of Earth's atmosphere.
As Mars approaches opposition, it begins a period of retrograde motion, which means it will appear to move backwards in a looping curve with respect to the background stars. This retrograde motion lasts for about 72 days, and Mars reaches its peak apparent brightness in the middle of this interval.
Moons
Mars has two relatively small (compared to Earth's) natural moons, Phobos (about in diameter) and Deimos (about in diameter), which orbit close to the planet. The origin of both moons is unclear, although a popular theory states that they were asteroids captured into Martian orbit.
Both satellites were discovered in 1877 by Asaph Hall and were named after the characters Phobos (the deity of panic and fear) and Deimos (the deity of terror and dread), twins from Greek mythology who accompanied their father Ares, god of war, into battle. Mars was the Roman equivalent to Ares. In modern Greek, the planet retains its ancient name Ares (Aris: Άρης).
From the surface of Mars, the motions of Phobos and Deimos appear different from that of the Earth's satellite, the Moon. Phobos rises in the west, sets in the east, and rises again in just 11 hours. Deimos, being only just outside synchronous orbitwhere the orbital period would match the planet's period of rotationrises as expected in the east, but slowly. Because the orbit of Phobos is below a synchronous altitude, tidal forces from Mars are gradually lowering its orbit. In about 50 million years, it could either crash into Mars's surface or break up into a ring structure around the planet.
The origin of the two satellites is not well understood. Their low albedo and carbonaceous chondrite composition have been regarded as similar to asteroids, supporting a capture theory. The unstable orbit of Phobos would seem to point toward a relatively recent capture. But both have circular orbits near the equator, which is unusual for captured objects, and the required capture dynamics are complex. Accretion early in the history of Mars is plausible, but would not account for a composition resembling asteroids rather than Mars itself, if that is confirmed. Mars may have yet-undiscovered moons, smaller than in diameter, and a dust ring is predicted to exist between Phobos and Deimos.
A third possibility for their origin as satellites of Mars is the involvement of a third body or a type of impact disruption. More-recent lines of evidence for Phobos having a highly porous interior, and suggesting a composition containing mainly phyllosilicates and other minerals known from Mars, point toward an origin of Phobos from material ejected by an impact on Mars that reaccreted in Martian orbit, similar to the prevailing theory for the origin of Earth's satellite. Although the visible and near-infrared (VNIR) spectra of the moons of Mars resemble those of outer-belt asteroids, the thermal infrared spectra of Phobos are reported to be inconsistent with chondrites of any class. It is also possible that Phobos and Deimos were fragments of an older moon, formed by debris from a large impact on Mars, and then destroyed by a more recent impact upon the satellite.
Human observations and exploration
The history of observations of Mars is marked by oppositions of Mars when the planet is closest to Earth and hence is most easily visible, which occur every couple of years. Even more notable are the perihelic oppositions of Mars, which are distinguished because Mars is close to perihelion, making it even closer to Earth.
Ancient and medieval observations
The ancient Sumerians named Mars Nergal, the god of war and plague. During Sumerian times, Nergal was a minor deity of little significance, but, during later times, his main cult center was the city of Nineveh. In Mesopotamian texts, Mars is referred to as the "star of judgement of the fate of the dead". The existence of Mars as a wandering object in the night sky was also recorded by the ancient Egyptian astronomers and, by 1534 BCE, they were familiar with the retrograde motion of the planet. By the period of the Neo-Babylonian Empire, the Babylonian astronomers were making regular records of the positions of the planets and systematic observations of their behavior. For Mars, they knew that the planet made 37 synodic periods, or 42 circuits of the zodiac, every 79 years. They invented arithmetic methods for making minor corrections to the predicted positions of the planets. In Ancient Greece, the planet was known as . Commonly, the Greek name for the planet now referred to as Mars, was Ares. It was the Romans who named the planet Mars, for their god of war, often represented by the sword and shield of the planet's namesake.
In the fourth century BCE, Aristotle noted that Mars disappeared behind the Moon during an occultation, indicating that the planet was farther away. Ptolemy, a Greek living in Alexandria, attempted to address the problem of the orbital motion of Mars. Ptolemy's model and his collective work on astronomy was presented in the multi-volume collection later called the Almagest (from the Arabic for "greatest"), which became the authoritative treatise on Western astronomy for the next fourteen centuries. Literature from ancient China confirms that Mars was known by Chinese astronomers by no later than the fourth century BCE. In the East Asian cultures, Mars is traditionally referred to as the "fire star" based on the Wuxing system.
During the seventeenth century A.D., Tycho Brahe measured the diurnal parallax of Mars that Johannes Kepler used to make a preliminary calculation of the relative distance to the planet. From Brahe's observations of Mars, Kepler deduced that the planet orbited the Sun not in a circle, but in an ellipse. Moreover, Kepler showed that Mars sped up as it approached the Sun and slowed down as it moved farther away, in a manner that later physicists would explain as a consequence of the conservation of angular momentum. When the telescope became available, the diurnal parallax of Mars was again measured in an effort to determine the Sun-Earth distance. This was first performed by Giovanni Domenico Cassini in 1672. The early parallax measurements were hampered by the quality of the instruments. The only occultation of Mars by Venus observed was that of 13 October 1590, seen by Michael Maestlin at Heidelberg. In 1610, Mars was viewed by Italian astronomer Galileo Galilei, who was first to see it via telescope. The first person to draw a map of Mars that displayed any terrain features was the Dutch astronomer Christiaan Huygens.
Martian "canals"
By the 19th century, the resolution of telescopes reached a level sufficient for surface features to be identified. On 5 September 1877, a perihelic opposition to Mars occurred. The Italian astronomer Giovanni Schiaparelli used a telescope in Milan to help produce the first detailed map of Mars. These maps notably contained features he called canali, which, with the possible exception of the natural canyon Valles Marineris, were later shown to be an optical illusion. These canali were supposedly long, straight lines on the surface of Mars, to which he gave names of famous rivers on Earth. His term, which means "channels" or "grooves", was popularly mistranslated in English as "canals".
Influenced by the observations, the orientalist Percival Lowell founded an observatory which had 30- and 45-centimetre (12- and 18-in) telescopes. The observatory was used for the exploration of Mars during the last good opportunity in 1894, and the following less favorable oppositions. He published several books on Mars and life on the planet, which had a great influence on the public. The canali were independently observed by other astronomers, like Henri Joseph Perrotin and Louis Thollon in Nice, using one of the largest telescopes of that time.
The seasonal changes (consisting of the diminishing of the polar caps and the dark areas formed during Martian summers) in combination with the canals led to speculation about life on Mars, and it was a long-held belief that Mars contained vast seas and vegetation. As bigger telescopes were used, fewer long, straight canali were observed. During observations in 1909 by Antoniadi with an telescope, irregular patterns were observed, but no canali were seen.
Robotic exploration
Dozens of crewless spacecraft, including orbiters, landers, and rovers, have been sent to Mars by the Soviet Union, the United States, Europe, India, the United Arab Emirates, and China to study the planet's surface, climate, and geology. NASA's Mariner 4 was the first spacecraft to visit Mars; launched on 28 November 1964, it made its closest approach to the planet on 15 July 1965. Mariner 4 detected the weak Martian radiation belt, measured at about 0.1% that of Earth, and captured the first images of another planet from deep space.
Once spacecraft visited the planet during NASA's Mariner missions in the 1960s and 1970s, many previous concepts of Mars were radically broken. After the results of the Viking life-detection experiments, the hypothesis of a dead planet was generally accepted. The data from Mariner 9 and Viking allowed better maps of Mars to be made, and the Mars Global Surveyor mission, which launched in 1996 and operated until late 2006, produced complete, extremely detailed maps of the Martian topography, magnetic field and surface minerals. These maps are available online at websites including Google Mars. Both the Mars Reconnaissance Orbiter and Mars Express continued exploring with new instruments and supporting lander missions. NASA provides two online tools: Mars Trek, which provides visualizations of the planet using data from 50 years of exploration, and Experience Curiosity, which simulates traveling on Mars in 3-D with Curiosity.
, Mars is host to ten functioning spacecraft. Eight are in orbit: 2001 Mars Odyssey, Mars Express, Mars Reconnaissance Orbiter, MAVEN, ExoMars Trace Gas Orbiter, the Hope orbiter, and the Tianwen-1 orbiter. Another two are on the surface: the Mars Science Laboratory Curiosity rover and the Perseverance rover.
Planned missions to Mars include:
NASA's EscaPADE spacecraft, planned to launch in 2025.
The Rosalind Franklin rover mission, designed to search for evidence of past life, which was intended to be launched in 2018 but has been repeatedly delayed, with a launch date pushed to 2028 at the earliest. The project was restarted in 2024 with additional funding.
A current concept for a joint NASA-ESA mission to return samples from Mars would launch in 2026.
China's Tianwen-3, a sample return mission, scheduled to launch in 2030.
, debris from these types of missions has reached over seven tons. Most of it consists of crashed and inactive spacecraft as well as discarded components.
In April 2024, NASA selected several companies to begin studies on providing commercial services to further enable robotic science on Mars. Key areas include establishing telecommunications, payload delivery and surface imaging.
Habitability and the search for life
During the late 19th century, it was widely accepted in the astronomical community that Mars had life-supporting qualities, including the presence of oxygen and water. However, in 1894 W. W. Campbell at Lick Observatory observed the planet and found that "if water vapor or oxygen occur in the atmosphere of Mars it is in quantities too small to be detected by spectroscopes then available". That observation contradicted many of the measurements of the time and was not widely accepted. Campbell and V. M. Slipher repeated the study in 1909 using better instruments, but with the same results. It was not until the findings were confirmed by W. S. Adams in 1925 that the myth of the Earth-like habitability of Mars was finally broken. However, even in the 1960s, articles were published on Martian biology, putting aside explanations other than life for the seasonal changes on Mars.
The current understanding of planetary habitabilitythe ability of a world to develop environmental conditions favorable to the emergence of lifefavors planets that have liquid water on their surface. Most often this requires the orbit of a planet to lie within the habitable zone, which for the Sun is estimated to extend from within the orbit of Earth to about that of Mars. During perihelion, Mars dips inside this region, but Mars's thin (low-pressure) atmosphere prevents liquid water from existing over large regions for extended periods. The past flow of liquid water demonstrates the planet's potential for habitability. Recent evidence has suggested that any water on the Martian surface may have been too salty and acidic to support regular terrestrial life.
The environmental conditions on Mars are a challenge to sustaining organic life: the planet has little heat transfer across its surface, it has poor insulation against bombardment by the solar wind due to the absence of a magnetosphere and has insufficient atmospheric pressure to retain water in a liquid form (water instead sublimes to a gaseous state). Mars is nearly, or perhaps totally, geologically dead; the end of volcanic activity has apparently stopped the recycling of chemicals and minerals between the surface and interior of the planet.
Evidence suggests that the planet was once significantly more habitable than it is today, but whether living organisms ever existed there remains unknown. The Viking probes of the mid-1970s carried experiments designed to detect microorganisms in Martian soil at their respective landing sites and had positive results, including a temporary increase in production on exposure to water and nutrients. This sign of life was later disputed by scientists, resulting in a continuing debate, with NASA scientist Gilbert Levin asserting that Viking may have found life. A 2014 analysis of Martian meteorite EETA79001 found chlorate, perchlorate, and nitrate ions in sufficiently high concentrations to suggest that they are widespread on Mars. UV and X-ray radiation would turn chlorate and perchlorate ions into other, highly reactive oxychlorines, indicating that any organic molecules would have to be buried under the surface to survive.
Small quantities of methane and formaldehyde detected by Mars orbiters are both claimed to be possible evidence for life, as these chemical compounds would quickly break down in the Martian atmosphere. Alternatively, these compounds may instead be replenished by volcanic or other geological means, such as serpentinite. Impact glass, formed by the impact of meteors, which on Earth can preserve signs of life, has also been found on the surface of the impact craters on Mars. Likewise, the glass in impact craters on Mars could have preserved signs of life, if life existed at the site.
The Cheyava Falls rock discovered on Mars in June 2024 has been designated by NASA as a "potential biosignature" and was core sampled by the Perseverance rover for possible return to Earth and further examination. Although highly intriguing, no definitive final determination on a biological or abiotic origin of this rock can be made with the data currently available.
Human mission proposals
Several plans for a human mission to Mars have been proposed throughout the 20th and 21st centuries, but none have come to fruition. The NASA Authorization Act of 2017 directed NASA to study the feasibility of a crewed Mars mission in the early 2030s; the resulting report eventually concluded that this would be unfeasible. In addition, in 2021, China was planning to send a crewed Mars mission in 2033. Privately held companies such as SpaceX have also proposed plans to send humans to Mars, with the eventual goal to settle on the planet. As of 2024, SpaceX has proceeded with the development of the Starship launch vehicle with the goal of Mars colonization. In plans shared with the company in April 2024, Elon Musk envisions the beginning of a Mars colony within the next twenty years. This enabled by the planned mass manufacturing of Starship and initially sustained by resupply from Earth, and in situ resource utilization on Mars, until the Mars colony reaches full self sustainability. Any future human mission to Mars will likely take place within the optimal Mars launch window, which occurs every 26 months. The moon Phobos has been proposed as an anchor point for a space elevator. Besides national space agencies and space companies, there are groups such as the Mars Society and The Planetary Society that advocates for human missions to Mars.
In culture
Mars is named after the Roman god of war (Greek Ares), but was also associated with the demi-god Heracles (Roman Hercules) by ancient Greek astronomers, as detailed by Aristotle. This association between Mars and war dates back at least to Babylonian astronomy, in which the planet was named for the god Nergal, deity of war and destruction. It persisted into modern times, as exemplified by Gustav Holst's orchestral suite The Planets, whose famous first movement labels Mars "the bringer of war". The planet's symbol, a circle with a spear pointing out to the upper right, is also used as a symbol for the male gender. The symbol dates from at least the 11th century, though a possible predecessor has been found in the Greek Oxyrhynchus Papyri.
The idea that Mars was populated by intelligent Martians became widespread in the late 19th century. Schiaparelli's "canali" observations combined with Percival Lowell's books on the subject put forward the standard notion of a planet that was a drying, cooling, dying world with ancient civilizations constructing irrigation works. Many other observations and proclamations by notable personalities added to what has been termed "Mars Fever". High-resolution mapping of the surface of Mars revealed no artifacts of habitation, but pseudoscientific speculation about intelligent life on Mars still continues. Reminiscent of the canali observations, these speculations are based on small scale features perceived in the spacecraft images, such as "pyramids" and the "Face on Mars". In his book Cosmos, planetary astronomer Carl Sagan wrote: "Mars has become a kind of mythic arena onto which we have projected our Earthly hopes and fears."
The depiction of Mars in fiction has been stimulated by its dramatic red color and by nineteenth-century scientific speculations that its surface conditions might support not just life but intelligent life. This gave way to many science fiction stories involving these concepts, such as H. G. Wells's The War of the Worlds, in which Martians seek to escape their dying planet by invading Earth; Ray Bradbury's The Martian Chronicles, in which human explorers accidentally destroy a Martian civilization; as well as Edgar Rice Burroughs's series Barsoom, C. S. Lewis's novel Out of the Silent Planet (1938), and a number of Robert A. Heinlein stories before the mid-sixties. Since then, depictions of Martians have also extended to animation. A comic figure of an intelligent Martian, Marvin the Martian, appeared in Haredevil Hare (1948) as a character in the Looney Tunes animated cartoons of Warner Brothers, and has continued as part of popular culture to the present. After the Mariner and Viking spacecraft had returned pictures of Mars as a lifeless and canal-less world, these ideas about Mars were abandoned; for many science-fiction authors, the new discoveries initially seemed like a constraint, but eventually the post-Viking knowledge of Mars became itself a source of inspiration for works like Kim Stanley Robinson's Mars trilogy.
| Physical sciences | Astronomy | null |
14644243 | https://en.wikipedia.org/wiki/Taxonomic%20rank | Taxonomic rank | In biology, taxonomic rank (which some authors prefer to call nomenclatural rank because ranking is part of nomenclature rather than taxonomy proper, according to some definitions of these terms) is the relative or absolute level of a group of organisms (a taxon) in a hierarchy that reflects evolutionary relationships. Thus, the most inclusive clades (such as Eukarya and Opisthokonta) have the highest ranks, whereas the least inclusive ones (such as Homo sapiens or Bufo bufo) have the lowest ranks. Ranks can be either relative and be denoted by an indented taxonomy in which the level of indentation reflects the rank, or absolute, in which various terms, such as species, genus, family, order, class, phylum, kingdom, and domain designate rank. This page emphasizes absolute ranks and the rank-based codes (the Zoological Code, the Botanical Code, the Code for Cultivated Plants, the Prokaryotic Code, and the Code for Viruses) require them. However, absolute ranks are not required in all nomenclatural systems for taxonomists; for instance, the PhyloCode, the code of phylogenetic nomenclature, does not require absolute ranks.
Taxa are hierarchical groups of organisms, and their ranks describes their position in this hierarchy. High-ranking taxa (e.g. those considered to be domains or kingdoms, for instance) include more sub-taxa than low-ranking taxa (e.g. those considered genera, species or subspecies). The rank of these taxa reflects inheritance of traits or molecular features from common ancestors. The name of any species and genus are basic; which means that to identify a particular organism, it is usually not necessary to specify names at ranks other than these first two, within a set of taxa covered by a given rank-based code. However, this is not true globally because most rank-based codes are independent from each other, so there are many inter-code homonyms (the same name used for different organisms, often for an animal and for a taxon covered by the botanical code). For this reason, attempts were made at creating a BioCode that would regulate all taxon names, but this attempt has so far failed because of firmly entrenched traditions in each community.
Consider a particular species, the red fox, Vulpes vulpes: in the context of the Zoological Code, the specific epithet vulpes (small v) identifies a particular species in the genus Vulpes (capital V) which comprises all the "true" foxes. Their close relatives are all in the family Canidae, which includes dogs, wolves, jackals, and all foxes; the next higher major taxon, Carnivora (considered an order), includes caniforms (bears, seals, weasels, skunks, raccoons and all those mentioned above), and feliforms (cats, civets, hyenas, mongooses). Carnivorans are one group of the hairy, warm-blooded, nursing members of the class Mammalia, which are classified among animals with notochords in the phylum Chordata, and with them among all animals in the kingdom Animalia. Finally, at the highest rank all of these are grouped together with all other organisms possessing cell nuclei in the domain Eukarya.
The International Code of Zoological Nomenclature defines rank as: "The level, for nomenclatural purposes, of a taxon in a taxonomic hierarchy (e.g. all families are for nomenclatural purposes at the same rank, which lies between superfamily and subfamily)." Note that the discussions on this page generally assume that taxa are clades (monophyletic groups of organisms), but this is required neither by the International Code of Zoological Nomenclature nor by the Botanical Code, and some experts on biological nomenclature do not think that this should be required, and in that case, the hierarchy of taxa (hence, their ranks) does not necessarily reflect the hierarchy of clades.
History
While older approaches to taxonomic classification were phenomenological, forming groups on the basis of similarities in appearance, organic structure and behavior, two important new methods developed in the second half of the 20th century changed drastically taxonomic practice. One is the advent of cladistics, which stemmed from the works of the German entomologist Willi Hennig. Cladistics is a method of classification of life forms according to the proportion of characteristics that they have in common (called synapomorphies). It is assumed that the higher the proportion of characteristics that two organisms share, the more recently they both came from a common ancestor. The second one is molecular systematics, based on genetic analysis, which can provide much additional data that prove especially useful when few phenotypic characters can resolve relationships, as, for instance, in many viruses, bacteria and archaea, or to resolve relationships between taxa that arose in a fast evolutionary radiation that occurred long ago, such as the main taxa of placental mammals.
Main ranks
In his landmark publications, such as the Systema Naturae, Carl Linnaeus used a ranking scale limited to kingdom, class, order, genus, species, and one rank below species. Today, the nomenclature is regulated by the nomenclature codes. There are seven main taxonomic ranks: kingdom, phylum or division, class, order, family, genus, and species. In addition, domain (proposed by Carl Woese) is now widely used as a fundamental rank, although it is not mentioned in any of the nomenclature codes, and is a synonym for dominion (), introduced by Moore in 1974.
A taxon is usually assigned a rank when it is given its formal name. The basic ranks are species and genus. When an organism is given a species name it is assigned to a genus, and the genus name is part of the species name.
The species name is also called a binomial, that is, a two-term name. For example, the zoological name for the human species is Homo sapiens. This is usually italicized in print or underlined when italics are not available. In this case, Homo is the generic name and it is capitalized; sapiens indicates the species and it is not capitalized. While not always used, some species include a subspecific epithet. For instance, modern humans are Homo sapiens sapiens, or H. sapiens sapiens.
In zoological nomenclature, higher taxon names are normally not italicized, but the Botanical Code, the Prokaryotic Code, the Code for Viruses, the draft BioCode and the PhyloCode all recommend italicizing all taxon names (of all ranks).
Ranks in zoology
There are rules applying to the following taxonomic ranks in the International Code of Zoological Nomenclature: superfamily, family, subfamily, tribe, subtribe, genus, subgenus, species, subspecies.
The International Code of Zoological Nomenclature divides names into "family-group names", "genus-group names" and "species-group names". The Code explicitly mentions the following ranks for these categories:
Family-groups
Superfamily (-oidea)
Family (-idae)
Subfamily (-inae)
Tribe (-ini)
Subtribe (-ina)
Genus-groups
Genus
Subgenus
Species-groups
Species
Subspecies
The rules in the Code apply to the ranks of superfamily to subspecies, and only to some extent to those above the rank of superfamily. Among "genus-group names" and "species-group names" no further ranks are officially allowed, which creates problems when naming taxa in these groups in speciose clades, such as Rana. Zoologists sometimes use additional terms such as species group, species subgroup, species complex and superspecies for convenience as extra, but unofficial, ranks between the subgenus and species levels in taxa with many species, e.g. the genus Drosophila. (Note the potentially confusing use of "species group" as both a category of ranks as well as an unofficial rank itself. For this reason, Alain Dubois has been using the alternative expressions "nominal-series", "family-series", "genus-series" and "species-series" (among others) at least since 2000.)
At higher ranks (family and above) a lower level may be denoted by adding the prefix "infra", meaning lower, to the rank. For example, infraorder (below suborder) or infrafamily (below subfamily).
Names of zoological taxa
A taxon above the rank of species has a scientific name in one part (a uninominal name).
A species has a name typically composed of two parts (a binomial name or binomen): generic name + specific name; for example Canis lupus. Sometimes the name of a subgenus (in parentheses) can be intercalated between the genus name and the specific epithet, which yields a trinomial name that should not be confused with that of a subspecies. An example is Lithobates (Aquarana) catesbeianus, which designates a species that belongs to the genus Lithobates and the subgenus Aquarana.
A subspecies has a name composed of three parts (a trinomial name or trinomen): generic name + specific name + subspecific name; for example Canis lupus italicus. As there is only one possible rank below that of species, no connecting term to indicate rank is needed or used.
Ranks in botany
Botanical ranks categorize organisms based (often) on their relationships (monophyly is not required by that clade, which does not even mention this word, nor that of "clade"). They start with Kingdom, then move to Division (or Phylum), Class, Order, Family, Genus, and Species. Taxa at each rank generally possess shared characteristics and evolutionary history. Understanding these ranks aids in taxonomy and studying biodiversity.
There are definitions of the following taxonomic categories in the International Code of Nomenclature for Cultivated Plants: cultivar group, cultivar, grex.
The rules in the ICN apply primarily to the ranks of family and below, and only to some extent to those above the rank of family.
Names of botanical taxa
Taxa at the rank of genus and above have a botanical name in one part (unitary name); those at the rank of species and above (but below genus) have a botanical name in two parts (binary name); all taxa below the rank of species have a botanical name in three parts (an infraspecific name). To indicate the rank of the infraspecific name, a "connecting term" is needed. Thus Poa secunda subsp. juncifolia, where "subsp". is an abbreviation for "subspecies", is the name of a subspecies of Poa secunda.
Hybrids can be specified either by a "hybrid formula" that specifies the parentage, or may be given a name. For hybrids receiving a hybrid name, the same ranks apply, prefixed with notho (Greek: 'bastard'), with nothogenus as the highest permitted rank.
Outdated names for botanical ranks
If a different term for the rank was used in an old publication, but the intention is clear, botanical nomenclature specifies certain substitutions:
If names were "intended as names of orders, but published with their rank denoted by a term such as": "cohors" [Latin for "cohort"; see also cohort study for the use of the term in ecology], "nixus", "alliance", or "Reihe" instead of "order" (Article 17.2), they are treated as names of orders.
"Family" is substituted for "order" (ordo) or "natural order" (ordo naturalis) under certain conditions where the modern meaning of "order" was not intended. (Article 18.2)
"Subfamily" is substituted for "suborder" (subordo) under certain conditions where the modern meaning of "suborder" was not intended. (Article 19.2)
In a publication prior to 1 January 1890, if only one infraspecific rank is used, it is considered to be that of variety. (Article 37.4) This commonly applies to publications that labelled infraspecific taxa with Greek letters, α, β, γ, ...
Examples
Classifications of five species follow: the fruit fly familiar in genetics laboratories (Drosophila melanogaster), humans (Homo sapiens), the peas used by Gregor Mendel in his discovery of genetics (Pisum sativum), the "fly agaric" mushroom Amanita muscaria, and the bacterium Escherichia coli. The eight major ranks are given in bold; a selection of minor ranks are given as well.
Table notes
In order to keep the table compact and avoid disputed technicalities, some common and uncommon intermediate ranks are omitted. For example, the mammals of Europe, Africa, and upper North America are in class Mammalia, legion Cladotheria, sublegion Zatheria, infralegion Tribosphenida, subclass Theria, clade Eutheria, clade Placentalia – but only Mammalia and Theria are in the table. Legitimate arguments might arise if the commonly used clades Eutheria and Placentalia were both included, over which is the rank "infraclass" and what the other's rank should be, or whether the two names are synonyms.
The ranks of higher taxa, especially intermediate ranks, are prone to revision as new information about relationships is discovered. For example, the flowering plants have been downgraded from a division (Magnoliophyta) to a subclass (Magnoliidae), and the superorder has become the rank that distinguishes the major groups of flowering plants. The traditional classification of primates (class Mammalia, subclass Theria, infraclass Eutheria, order Primates) has been modified by new classifications such as McKenna and Bell (class Mammalia, subclass Theriformes, infraclass Holotheria) with Theria and Eutheria assigned lower ranks between infraclass and the order Primates. These differences arise because there are few available ranks and many branching points in the fossil record.
Within species further units may be recognised. Animals may be classified into subspecies (for example, Homo sapiens sapiens, modern humans) or morphs (for example Corvus corax varius morpha leucophaeus, the pied raven). Plants may be classified into subspecies (for example, Pisum sativum subsp. sativum, the garden pea) or varieties (for example, Pisum sativum var. macrocarpon, snow pea), with cultivated plants getting a cultivar name (for example, Pisum sativum var. macrocarpon 'Snowbird'). Bacteria may be classified by strains (for example Escherichia coli O157:H7, a strain that can cause food poisoning).
Terminations of names
Taxa above the genus level are often given names based on the type genus, with a standard termination. The terminations used in forming these names depend on the kingdom (and sometimes the phylum and class) as set out in the table below.
Pronunciations given are the most Anglicized. More Latinate pronunciations are also common, particularly rather than for stressed a.
Table notes
In botany and mycology names at the rank of family and below are based on the name of a genus, sometimes called the type genus of that taxon, with a standard ending. For example, the rose family, Rosaceae, is named after the genus Rosa, with the standard ending "-aceae" for a family. Names above the rank of family are also formed from a generic name, or are descriptive (like Gymnospermae or Fungi).
For animals, there are standard suffixes for taxa only up to the rank of superfamily. Uniform suffix has been suggested (but not recommended) in AAAS as -ida for orders, for example; protozoologists seem to adopt this system. Many metazoan (higher animals) orders also have such suffix, e.g. Hyolithida and Nectaspida (Naraoiida).
Forming a name based on a generic name may be not straightforward. For example, the has the genitive , thus the genus Homo (human) is in the Hominidae, not "Homidae".
The ranks of epifamily, infrafamily and infratribe (in animals) are used where the complexities of phyletic branching require finer-than-usual distinctions. Although they fall below the rank of superfamily, they are not regulated under the International Code of Zoological Nomenclature and hence do not have formal standard endings. The suffixes listed here are regular, but informal.
In virology, the formal endings for taxa of viroids, of satellite nucleic acids, and of viriforms are similar to viruses, only -vir- is replaced by -viroid-, -satellit- and -viriform-. The extra levels of realm and subrealm end with -viria and -vira respectively.
All ranks
There is an indeterminate number of ranks, as a taxonomist may invent a new rank at will, at any time, if they feel this is necessary. In doing so, there are some restrictions, which will vary with the nomenclature code that applies.
The following is an artificial synthesis, solely for purposes of demonstration of absolute rank (but see notes), from most general to most specific:
Superdomain
Domain or Empire
Subdomain (biology)
Realm (in virology)
Subrealm (in virology)
Hyperkingdom
Superkingdom
Kingdom
Subkingdom
Infrakingdom
Parvkingdom
Superphylum, or superdivision (in botany)
Phylum, or division (in botany)
Subphylum, or subdivision (in botany)
Infraphylum, or infradivision (in botany)
Microphylum
Superclass
Class
Subclass
Infraclass
Subterclass
Parvclass
Superdivision (in zoology)
Division (in zoology)
Subdivision (in zoology)
Infradivision (in zoology)
Superlegion (in zoology)
Legion (in zoology)
Sublegion (in zoology)
Infralegion (in zoology)
Supercohort (in zoology)
Cohort (in zoology)
Subcohort (in zoology)
Infracohort (in zoology)
Gigaorder (in zoology)
Magnorder or megaorder (in zoology)
Grandorder or capaxorder (in zoology)
Mirorder or hyperorder (in zoology)
Superorder
Series (in ichthyology)
Order
Parvorder (position in some zoological classifications)
Nanorder (in zoology)
Hypoorder (in zoology)
Minorder (in zoology)
Suborder
Infraorder
Parvorder (usual position), or microorder (in zoology)
Section (in zoology)
Subsection (in zoology)
Gigafamily (in zoology)
Megafamily (in zoology)
Grandfamily (in zoology)
Hyperfamily (in zoology)
Superfamily
Epifamily (in zoology)
Series (for Lepidoptera)
Group (for Lepidoptera)
Family
Subfamily
Infrafamily
Supertribe
Tribe
Subtribe
Infratribe
Supergenus
Genus
Subgenus
Section (in botany)
Subsection (in botany)
Series (in botany)
Subseries (in botany)
Species complex
Species
Subspecies, or forma specialis (for fungi), or pathovar (for bacteria))
Variety or varietas (in botany); or form or morph (in zoology) or aberration (in lepidopterology)
Subvariety (in botany)
Form or forma (in botany)
Subform or subforma (in botany)
Significance and problems
Ranks are assigned based on subjective dissimilarity, and do not fully reflect the gradational nature of variation within nature. These problems were already identified by Willi Hennig, who advocated dropping them in 1969, and this position gathered support from Graham C. D. Griffiths only a few years later. In fact, these ranks were proposed in a fixist context and the advent of evolution sapped the foundations of this system, as was recognised long ago; the introduction of The Code of Nomenclature and Check-list of North American Birds Adopted by the American Ornithologists' Union published in 1886 states "No one appears to have suspected, in 1842 [when the Strickland code was drafted], that the Linnaean system was not the permanent heritage of science, or that in a few years a theory of evolution was to sap its very foundations, by radically changing men's conceptions of those things to which names were to be furnished." Such ranks are used simply because they are required by the rank-based codes; because of this, some systematists prefer to call them nomenclatural ranks. In most cases, higher taxonomic groupings arise further back in time, simply because the most inclusive taxa necessarily appeared first. Furthermore, the diversity in some major taxa (such as vertebrates and angiosperms) is better known than that of others (such as fungi, arthropods and nematodes) not because they are more diverse than other taxa, but because they are more easily sampled and studied than other taxa, or because they attract more interest and funding for research.
Of these many ranks, many systematists consider that the most basic (or important) is the species, but this opinion is not universally shared. Thus, species are not necessarily more sharply defined than taxa at any other rank, and in fact, given the phenotypic gaps created by extinction, in practice, the reverse is often the case. Ideally, a taxon is intended to represent a clade, that is, the phylogeny of the organisms under discussion, but this is not a requirement of the zoological and botanical codes.
A classification in which all taxa have formal ranks cannot adequately reflect knowledge about phylogeny. Since taxon names are dependent on ranks in rank-based (Linnaean) nomenclature, taxa without ranks cannot be given names. Alternative approaches, such as phylogenetic nomenclature, as implemented under the PhyloCode and supported by the International Society for Phylogenetic Nomenclature, or using circumscriptional names, avoid this problem. The theoretical difficulty with superimposing taxonomic ranks over evolutionary trees is manifested as the boundary paradox which may be illustrated by Darwinian evolutionary models.
There are no rules for how many species should make a genus, a family, or any other higher taxon (that is, a taxon in a category above the species level). It should be a natural group (that is, non-artificial, non-polyphyletic), as judged by a biologist, using all the information available to them. Equally ranked higher taxa in different phyla are not necessarily equivalent in terms of time of origin, phenotypic distinctiveness or number of lower-ranking included taxa (e.g., it is incorrect to assume that families of insects are in some way evolutionarily comparable to families of mollusks). Of all criteria that have been advocated to rank taxa, age of origin has been the most frequently advocated. Willi Hennig proposed it in 1966, but he concluded in 1969 that this system was unworkable and suggested dropping absolute ranks. However, the idea of ranking taxa using the age of origin (either as the sole criterion, or as one of the main ones) persists under the name of time banding, and is still advocated by several authors. For animals, at least the phylum rank is usually associated with a certain body plan, which is also, however, an arbitrary criterion.
Enigmatic taxa
Enigmatic taxa are taxonomic groups whose broader relationships are unknown or undefined.
Mnemonic
There are several acronyms intended to help memorise the taxonomic hierarchy, such as "King Phillip came over for great spaghetti".
| Biology and health sciences | Taxonomic rank | Biology |
169358 | https://en.wikipedia.org/wiki/Foundations%20of%20mathematics | Foundations of mathematics | Foundations of mathematics are the logical and mathematical framework that allows the development of mathematics without generating self-contradictory theories, and, in particular, to have reliable concepts of theorems, proofs, algorithms, etc. This may also include the philosophical study of the relation of this framework with reality.
The term "foundations of mathematics" was not coined before the end of the 19th century, although foundations were first established by the ancient Greek philosophers under the name of Aristotle's logic and systematically applied in Euclid's Elements. A mathematical assertion is considered as truth only if it is a theorem that is proved from true premises by means of a sequence of syllogisms (inference rules), the premises being either already proved theorems or self-evident assertions called axioms or postulates.
These foundations were tacitly assumed to be definitive until the introduction of infinitesimal calculus by Isaac Newton and Gottfried Wilhelm Leibniz in the 17th century. This new area of mathematics involved new methods of reasoning and new basic concepts (continuous functions, derivatives, limits) that were not well founded, but had astonishing consequences, such as the deduction from Newton's law of gravitation that the orbits of the planets are ellipses.
During the 19th century, progress was made towards elaborating precise definitions of the basic concepts of infinitesimal calculus, notably the natural and real numbers. This led, near the end of the 19th century, to a series of seemingly paradoxical mathematical results that challenged the general confidence in the reliability and truth of mathematical results. This has been called the foundational crisis of mathematics.
The resolution of this crisis involved the rise of a new mathematical discipline called mathematical logic that includes set theory, model theory, proof theory, computability and computational complexity theory, and more recently, parts of computer science. Subsequent discoveries in the 20th century then stabilized the foundations of mathematics into a coherent framework valid for all mathematics. This framework is based on a systematic use of axiomatic method and on set theory, specifically ZFC, the Zermelo–Fraenkel set theory with the axiom of choice.
It results from this that the basic mathematical concepts, such as numbers, points, lines, and geometrical spaces are not defined as abstractions from reality but from basic properties (axioms). Their adequation with their physical origins does not belong to mathematics anymore, although their relation with reality is still used for guiding mathematical intuition: physical reality is still used by mathematicians to choose axioms, find which theorems are interesting to prove, and obtain indications of possible proofs.
Ancient Greece
Most civilisations developed some mathematics, mainly for practical purposes, such as counting (merchants), surveying (delimitation of fields), prosody, astronomy, and astrology. It seems that ancient Greek philosophers were the first to study the nature of mathematics and its relation with the real world.
Zeno of Elea (490 c. 430 BC) produced several paradoxes he used to support his thesis that movement does not exist. These paradoxes involve mathematical infinity, a concept that was outside the mathematical foundations of that time and was not well understood before the end of the 19th century.
The Pythagorean school of mathematics originally insisted that the only numbers are natural numbers and ratios of natural numbers. The discovery (around 5th century BC) that the ratio of the diagonal of a square to its side is not the ratio of two natural numbers was a shock to them which they only reluctantly accepted. A testimony of this is the modern terminology of irrational number for referring to a number that is not the quotient of two integers, since "irrational" means originally "not reasonable" or "not accessible with reason".
The fact that length ratios are not represented by rational numbers was resolved by Eudoxus of Cnidus (408–355 BC), a student of Plato, who reduced the comparison of two irrational ratios to comparisons of integer multiples of the magnitudes involved. His method anticipated that of Dedekind cuts in the modern definition of real numbers by Richard Dedekind (1831–1916); see .
In the Posterior Analytics, Aristotle (384–322 BC) laid down the logic for organizing a field of knowledge by means of primitive concepts, axioms, postulates, definitions, and theorems. Aristotle took a majority of his examples for this from arithmetic and from geometry, and his logic served as the foundation of mathematics for centuries. This method resembles the modern axiomatic method but with a big philosophical difference: axioms and postulates were supposed to be true, being either self-evident or resulting from experiments, while no other truth than the correctness of the proof is involved in the axiomatic method. So, for Aristotle, a proved theorem is true, while in the axiomatic methods, the proof says only that the axioms imply the statement of the theorem.
Aristotle's logic reached its high point with Euclid's Elements (300 BC), a treatise on mathematics structured with very high standards of rigor: Euclid justifies each proposition by a demonstration in the form of chains of syllogisms (though they do not always conform strictly to Aristotelian templates).
Aristotle's syllogistic logic, together with its exemplification by Euclid's Elements, are recognized as scientific achievements of ancient Greece, and remained as the foundations of mathematics for centuries.
Before infinitesimal calculus
During Middle Ages, Euclid's Elements stood as a perfectly solid foundation for mathematics, and philosophy of mathematics concentrated on the ontological status of mathematical concepts; the question was whether they exist independently of perception (realism) or within the mind only (conceptualism); or even whether they are simply names of collection of individual objects (nominalism).
In Elements, the only numbers that are considered are natural numbers and ratios of lengths. This geometrical view of non-integer numbers remained dominant until the end of Middle Ages, although the rise of algebra led to consider them independently from geometry, which implies implicitly that there are foundational primitives of mathematics. For example, the transformations of equations introduced by Al-Khwarizmi and the cubic and quartic formulas discovered in the 16th century result from algebraic manipulations that have no geometric counterpart.
Nevertheless, this did not challenge the classical foundations of mathematics since all properties of numbers that were used can be deduced from their geometrical definition.
In 1637, René Descartes published La Géométrie, in which he showed that geometry can be reduced to algebra by means coordinates, which are numbers determining the position of a point. This gives to the numbers that he called real numbers a more foundational role (before him, numbers were defined as the ratio of two lengths). Descartes' book became famous after 1649 and paved the way to infinitesimal calculus.
Infinitesimal calculus
Isaac Newton (1642–1727) in England and Leibniz (1646–1716) in Germany independently developed the infinitesimal calculus for dealing with mobile points (such as planets in the sky) and variable quantities.
This needed the introduction of new concepts such as continuous functions, derivatives and limits. For dealing with these concepts in a logical way, they were defined in terms of infinitesimals that are hypothetical numbers that are infinitely close to zero. The strong implications of infinitesimal calculus on foundations of mathematics is illustrated by a pamphlet of the Protestant philosopher George Berkeley (1685–1753), who wrote "[Infinitesimals] are neither finite quantities, nor quantities infinitely small, nor yet nothing. May we not call them the ghosts of departed quantities?".
Also, a lack of rigor has been frequently invoked, because infinitesimals and the associated concepts were not formally defined (lines and planes were not formally defined either, but people were more accustomed to them). Real numbers, continuous functions, derivatives were not formally defined before the 19th century, as well as Euclidean geometry. It is only in the 20th century that a formal definition of infinitesimals has been given, with the proof that the whole infinitesimal can be deduced from them.
Despite its lack of firm logical foundations, infinitesimal calculus was quickly adopted by mathematicians, and validated by its numerous applications; in particular the fact that the planet trajectories can be deduced from the Newton's law of gravitation.
19th century
In the 19th century, mathematics developed quickly in many directions. Several of the problems that were considered led to questions on the foundations of mathematics. Frequently, the proposed solutions led to further questions that were often simultaneously of philosophical and mathematical nature. All these questions led, at the end of the 19th century and the beginning of the 20th century, to debates which have been called the foundational crisis of mathematics. The following subsections describe the main such foundational problems revealed during the 19th century.
Real analysis
Cauchy (1789–1857) started the project of giving rigorous bases to infinitesimal calculus. In particular, he rejected the heuristic principle that he called the generality of algebra, which consisted to apply properties of algebraic operations to infinite sequences without proper proofs. In his Cours d'Analyse (1821), he considered very small quantities, which could presently be called "sufficiently small quantities"; that is, a sentence such that "if is very small must be understood as "there is a (sufficiently large) natural number such that ". In the proofs he used this in a way that predated the modern (ε, δ)-definition of limit.
The modern (ε, δ)-definition of limits and continuous functions was first developed by Bolzano in 1817, but remained relatively unknown, and Cauchy probably did know Bolzano's work.
Karl Weierstrass (1815–1897) formalized and popularized the (ε, δ)-definition of limits, and discovered some pathological functions that seemed paradoxical at this time, such as continuous, nowhere-differentiable functions. Indeed, such functions contradict previous conceptions of a function as a rule for computation or a smooth graph.
At this point, the program of arithmetization of analysis (reduction of mathematical analysis to arithmetic and algebraic operations) advocated by Weierstrass was essentially completed, except for two points.
Firstly, a formal definition of real numbers was still lacking. Indeed, beginning with Richard Dedekind in 1858, several mathematicians worked on the definition of the real numbers, including Hermann Hankel, Charles Méray, and Eduard Heine, but this is only in 1872 that two independent complete definitions of real numbers were published: one by Dedekind, by means of Dedekind cuts; the other one by Georg Cantor as equivalence classes of Cauchy sequences.
Several problems were left open by these definitions, which contributed to the foundational crisis of mathematics. Firstly both definitions suppose that rational numbers and thus natural numbers are rigorously defined; this was done a few years later with Peano axioms. Secondly, both definitions involve infinite sets (Dedekind cuts and sets of the elements of a Cauchy sequence), and Cantor's set theory was published several years later.
The third problem is more subtle: and is related to the foundations of logic: classical logic is a first order logic; that is, quantifiers apply to variables representing individual elements, not to variables representing (infinite) sets of elements. The basic property of the completeness of the real numbers that is required for defining and using real numbers involves a quantification on infinite sets. Indeed, this property may be expressed either as for every infinite sequence of real numbers, if it is a Cauchy sequence, it has a limit that is a real number, or as every subset of the real numbers that is bounded has a least upper bound that is a real number. This need of quantification over infinite sets is one of the motivation of the development of higher-order logics during the first half of the 20th century.
Non-Euclidean geometries
Before the 19th century, there were many failed attempts to derive the parallel postulate from other axioms of geometry. In an attempt to prove that its negation leads to a contradiction, Johann Heinrich Lambert (1728–1777) started to build hyperbolic geometry and introduced the hyperbolic functions and computed the area of a hyperbolic triangle (where the sum of angles is less than 180°).
Continuing the construction of this new geometry, several mathematicians proved independently that if it is inconsistent, then Euclidean geometry is also inconsistent and thus that the parallel postulate cannot be proved. This was proved by Nikolai Lobachevsky in 1826, János Bolyai (1802–1860) in 1832 and Carl Friedrich Gauss (unpublished).
Later in the 19th century, the German mathematician Bernhard Riemann developed Elliptic geometry, another non-Euclidean geometry where no parallel can be found and the sum of angles in a triangle is more than 180°. It was proved consistent by defining points as pairs of antipodal points on a sphere (or hypersphere), and lines as great circles on the sphere.
These proofs of unprovability of the parallel postulate lead to several philosophical problems, the main one being that before this discovery, the parallel postulate and all its consequences were considered as true. So, the non-Euclidean geometries challenged the concept of mathematical truth.
Synthetic vs. analytic geometry
Since the introduction of analytic geometry by René Descartes in the 17th century, there were two approaches to geometry, the old one called synthetic geometry, and the new one, where everything is specified in terms of real numbers called coordinates.
Mathematicians did not worry much about the contradiction between these two approaches before the mid-nineteenth century, where there was "an acrimonious controversy between the proponents of synthetic and analytic methods in projective geometry, the two sides accusing each other of mixing projective and metric concepts". Indeed, there is no concept of distance in a projective space, and the cross-ratio, which is a number, is a basic concept of synthetic projective geometry.
Karl von Staudt developed a purely geometric approach to this problem by introducing "throws" that form what is presently called a field, in which the cross ratio can be expressed.
Apparently, the problem of the equivalence between analytic and synthetic approach was completely solved only with Emil Artin's book Geometric Algebra published in 1957. It was well known that, given a field , one may define affine and projective spaces over in terms of -vector spaces. In these spaces, the Pappus hexagon theorem holds. Conversely, if the Pappus hexagon theorem is included in the axioms of a plane geometry, then one can define a field such that the geometry is the same as the affine or projective geometry over .
Natural numbers
The work of making rigorous real analysis and the definition of real numbers, consisted of reducing everything to rational numbers and thus to natural numbers, since positive rational numbers are fractions of natural numbers. There was therefore a need of a formal definition of natural numbers, which imply as axiomatic theory of arithmetic. This was started with Charles Sanders Peirce in 1881 and Richard Dedekind in 1888, who defined a natural numbers as the cardinality of a finite set.. However, this involves set theory, which was not formalized at this time.
Giuseppe Peano provided in 1888 a complete axiomatisation based on the ordinal property of the natural numbers. The last Peano's axiom is the only one that induces logical difficulties, as it begin with either "if is a set then" or "if is a predicate then". So, Peano's axioms induce a quantification on infinite sets, and this means that Peano arithmetic is what is presently called a Second-order logic.
This was not well understood at that times, but the fact that infinity occurred in the definition of the natural numbers was a problem for many mathematicians of this time. For example, Henri Poincaré stated that axioms can only be demonstrated in their finite application, and concluded that it is "the power of the mind" which allows conceiving of the indefinite repetition of the same act. This applies in particular to the use of the last Peano axiom for showing that the successor function generates all natural numbers. Also, Leopold Kronecker said "God made the integers, all else is the work of man". This may be interpreted as "the integers cannot be mathematically defined".
Infinite sets
Before the second half of the 19th century, infinity was a philosophical concept that did not belong to mathematics. However, with the rise of infinitesimal calculus, mathematicians became to be accustomed to infinity, mainly through potential infinity, that is, as the result of an endless process, such as the definition of an infinite sequence, an infinite series or a limit. The possibility of an actual infinity was the subject of many philosophical disputes.
Sets, and more specially infinite sets were not considered as a mathematical concept; in particular, there was no fixed term for them. A dramatic change arose with the work of Georg Cantor who was the first mathematician to systematically study infinite sets. In particular, he introduced cardinal numbers that measure the size of infinite sets, and ordinal numbers that, roughly speaking, allow one to continue to count after having reach infinity. One of his major results is the discovery that there are strictly more real numbers than natural numbers (the cardinal of the continuum of the real numbers is greater than that of the natural numbers).
These results were rejected by many mathematicians and philosophers, and led to debates that are a part of the foundational crisis of mathematics.
The crisis was amplified with the Russel's paradox that asserts that the phrase "the set of all sets" is self-contradictory. This condradiction introduced a doubt on the consistency of all mathematics.
With the introduction of the Zermelo–Fraenkel set theory () and its adoption by the mathematical community, the doubt about the consistency was essentially removed, although consistency of set theory cannot be proved because of Gödel's incompleteness theorem.
Mathematical logic
In 1847, De Morgan published his laws and George Boole devised an algebra, now called Boolean algebra, that allows expressing Aristotle's logic in terms of formulas and algebraic operations. Boolean algebra is the starting point of mathematization logic and the basis of propositional calculus
Independently, in the 1870's, Charles Sanders Peirce and Gottlob Frege extended propositional calculus by introducing quantifiers, for building predicate logic.
Frege pointed out three desired properties of a logical theory:consistency (impossibility of proving contradictory statements), completeness (any statement is either provable or refutable; that is, its negation is provable), and decidability (there is a decision procedure to test every statement).
By near the turn of the century, Bertrand Russell popularized Frege's work and discovered Russel's paradox which implies that the phrase "the set of all sets" is self-contradictory. This paradox seemed to make the whole mathematics inconsistent and is one of the major causes of the foundational crisis of mathematics.
Foundational crisis
The foundational crisis of mathematics
arose at the end of the 19th century and the beginning of the 20th century with the discovery of several paradoxes or counter-intuitive results.
The first one was the proof that the parallel postulate cannot be proved. This results from a construction of a non-Euclidean geometry inside Euclidean geometry, whose inconsistency would imply the inconsistency of Euclidean geometry. A well known paradox is Russell's paradox, which shows that the phrase "the set of all sets that do not contain themselves" is self-contradictory. Other philosophical problems were the proof of the existence of mathematical objects that cannot be computed or explicitly described, and the proof of the existence of theorems of arithmetic that cannot be proved with Peano arithmetic.
Several schools of philosophy of mathematics were challenged with these problems in the 20th century, and are described below.
These problems were also studied by mathematicians, and this led to establish mathematical logic as a new area of mathematics, consisting of providing mathematical definitions to logics (sets of inference rules), mathematical and logical theories, theorems, and proofs, and of using mathematical methods to prove theorems about these concepts.
This led to unexpected results, such as Gödel's incompleteness theorems, which, roughly speaking, assert that, if a theory contains the standard arithmetic, it cannot be used to prove that it itself is not self-contradictory; and, if it is not self-contradictory, there are theorems that cannot be proved inside the theory, but are nevertheless true in some technical sense.
Zermelo–Fraenkel set theory with the axiom of choice (ZFC) is a logical theory established by Ernst Zermelo and Abraham Fraenkel. It became the standard foundation of modern mathematics, and, unless the contrary is explicitly specified, it is used in all modern mathematical texts, generally implicitly.
Simultaneously, the axiomatic method became a de facto standard: the proof of a theorem must result from explicit axioms and previously proved theorems by the application of clearly defined inference rules. The axioms need not correspond to some reality. Nevertheless, it is an open philosophical problem to explain why the axiom systems that lead to rich and useful theories are those resulting from abstraction from the physical reality or other mathematical theory.
In summary, the foundational crisis is essentially resolved, and this opens new philosophical problems. In particular, it cannot be proved that the new foundation (ZFC) is not self-contradictory. It is a general consensus that, if this would happen, the problem could be solved by a mild modification of ZFC.
Philosophical views
When the foundational crisis arose, there was much debate among mathematicians and logicians about what should be done for restoring confidence in mathematics. This involved philosophical questions about mathematical truth, the relationship of mathematics with reality, the reality of mathematical objects, and the nature of mathematics.
For the problem of foundations, there was two main options for trying to avoid paradoxes. The first one led to intuitionism and constructivism, and consisted to restrict the logical rules for remaining closer to intuition, while the second, which has been called formalism, considers that a theorem is true if it can be deduced from axioms by applying inference rules (formal proof), and that no "trueness" of the axioms is needed for the validity of a theorem.
Formalism
It has been claimed that formalists, such as David Hilbert (1862–1943), hold that mathematics is only a language and a series of games. Hilbert insisted that formalism, called "formula game" by him, is a fundamental part of mathematics, but that mathematics must not be reduced to formalism. Indeed, he used the words "formula game" in his 1927 response to L. E. J. Brouwer's criticisms:
Thus Hilbert is insisting that mathematics is not an arbitrary game with arbitrary rules; rather it must agree with how our thinking, and then our speaking and writing, proceeds.
The foundational philosophy of formalism, as exemplified by David Hilbert, is a response to the paradoxes of set theory, and is based on formal logic. Virtually all mathematical theorems today can be formulated as theorems of set theory. The truth of a mathematical statement, in this view, is represented by the fact that the statement can be derived from the axioms of set theory using the rules of formal logic.
Merely the use of formalism alone does not explain several issues: why we should use the axioms we do and not some others, why we should employ the logical rules we do and not some others, why "true" mathematical statements (e.g., the laws of arithmetic) appear to be true, and so on. Hermann Weyl posed these very questions to Hilbert:
In some cases these questions may be sufficiently answered through the study of formal theories, in disciplines such as reverse mathematics and computational complexity theory. As noted by Weyl, formal logical systems also run the risk of inconsistency; in Peano arithmetic, this arguably has already been settled with several proofs of consistency, but there is debate over whether or not they are sufficiently finitary to be meaningful. Gödel's second incompleteness theorem establishes that logical systems of arithmetic can never contain a valid proof of their own consistency. What Hilbert wanted to do was prove a logical system S was consistent, based on principles P that only made up a small part of S. But Gödel proved that the principles P could not even prove P to be consistent, let alone S.
Intuitionism
Intuitionists, such as L. E. J. Brouwer (1882–1966), hold that mathematics is a creation of the human mind. Numbers, like fairy tale characters, are merely mental entities, which would not exist if there were never any human minds to think about them.
The foundational philosophy of intuitionism or constructivism, as exemplified in the extreme by Brouwer and Stephen Kleene, requires proofs to be "constructive" in nature the existence of an object must be demonstrated rather than inferred from a demonstration of the impossibility of its non-existence. For example, as a consequence of this the form of proof known as reductio ad absurdum is suspect.
Some modern theories in the philosophy of mathematics deny the existence of foundations in the original sense. Some theories tend to focus on mathematical practice, and aim to describe and analyze the actual working of mathematicians as a social group. Others try to create a cognitive science of mathematics, focusing on human cognition as the origin of the reliability of mathematics when applied to the real world. These theories would propose to find foundations only in human thought, not in any objective outside construct. The matter remains controversial.
Logicism
Logicism is a school of thought, and research programme, in the philosophy of mathematics, based on the thesis that mathematics is an extension of logic or that some or all mathematics may be derived in a suitable formal system whose axioms and rules of inference are 'logical' in nature. Bertrand Russell and Alfred North Whitehead championed this theory initiated by Gottlob Frege and influenced by Richard Dedekind.
Set-theoretic Platonism
Many researchers in axiomatic set theory have subscribed to what is known as set-theoretic Platonism, exemplified by Kurt Gödel.
Several set theorists followed this approach and actively searched for axioms that may be considered as true for heuristic reasons and that would decide the continuum hypothesis. Many large cardinal axioms were studied, but the hypothesis always remained independent from them and it is now considered unlikely that CH can be resolved by a new large cardinal axiom. Other types of axioms were considered, but none of them has reached consensus on the continuum hypothesis yet. Recent work by Hamkins proposes a more flexible alternative: a set-theoretic multiverse allowing free passage between set-theoretic universes that satisfy the continuum hypothesis and other universes that do not.
Indispensability argument for realism
This argument by Willard Quine and Hilary Putnam says (in Putnam's shorter words),
However, Putnam was not a Platonist.
Rough-and-ready realism
Few mathematicians are typically concerned on a daily, working basis over logicism, formalism or any other philosophical position. Instead, their primary concern is that the mathematical enterprise as a whole always remains productive. Typically, they see this as ensured by remaining open-minded, practical and busy; as potentially threatened by becoming overly-ideological, fanatically reductionistic or lazy.
Such a view has also been expressed by some well-known physicists.
For example, the Physics Nobel Prize laureate Richard Feynman said
And Steven Weinberg:
Weinberg believed that any undecidability in mathematics, such as the continuum hypothesis, could be potentially resolved despite the incompleteness theorem, by finding suitable further axioms to add to set theory.
Philosophical consequences of Gödel's completeness theorem
Gödel's completeness theorem establishes an equivalence in first-order logic between the formal provability of a formula and its truth in all possible models. Precisely, for any consistent first-order theory it gives an "explicit construction" of a model described by the theory; this model will be countable if the language of the theory is countable. However this "explicit construction" is not algorithmic. It is based on an iterative process of completion of the theory, where each step of the iteration consists in adding a formula to the axioms if it keeps the theory consistent; but this consistency question is only semi-decidable (an algorithm is available to find any contradiction but if there is none this consistency fact can remain unprovable).
More paradoxes
The following lists some notable results in metamathematics. Zermelo–Fraenkel set theory is the most widely studied axiomatization of set theory. It is abbreviated ZFC when it includes the axiom of choice and ZF when the axiom of choice is excluded.
1920: Thoralf Skolem corrected Leopold Löwenheim's proof of what is now called the downward Löwenheim–Skolem theorem, leading to Skolem's paradox discussed in 1922, namely the existence of countable models of ZF, making infinite cardinalities a relative property.
1922: Proof by Abraham Fraenkel that the axiom of choice cannot be proved from the axioms of Zermelo set theory with urelements.
1931: Publication of Gödel's incompleteness theorems, showing that essential aspects of Hilbert's program could not be attained. It showed how to construct, for any sufficiently powerful and consistent recursively axiomatizable system such as necessary to axiomatize the elementary theory of arithmetic on the (infinite) set of natural numbers a statement that formally expresses its own unprovability, which he then proved equivalent to the claim of consistency of the theory; so that (assuming the consistency as true), the system is not powerful enough for proving its own consistency, let alone that a simpler system could do the job. It thus became clear that the notion of mathematical truth cannot be completely determined and reduced to a purely formal system as envisaged in Hilbert's program. This dealt a final blow to the heart of Hilbert's program, the hope that consistency could be established by finitistic means (it was never made clear exactly what axioms were the "finitistic" ones, but whatever axiomatic system was being referred to, it was a 'weaker' system than the system whose consistency it was supposed to prove).
1936: Alfred Tarski proved his truth undefinability theorem.
1936: Alan Turing proved that a general algorithm to solve the halting problem for all possible program-input pairs cannot exist.
1938: Gödel proved the consistency of the axiom of choice and of the generalized continuum hypothesis.
1936–1937: Alonzo Church and Alan Turing, respectively, published independent papers showing that a general solution to the Entscheidungsproblem is impossible: the universal validity of statements in first-order logic is not decidable (it is only semi-decidable as given by the completeness theorem).
1955: Pyotr Novikov showed that there exists a finitely presented group G such that the word problem for G is undecidable.
1963: Paul Cohen showed that the Continuum Hypothesis is unprovable from ZFC. Cohen's proof developed the method of forcing, which is now an important tool for establishing independence results in set theory.
1964: Inspired by the fundamental randomness in physics, Gregory Chaitin starts publishing results on algorithmic information theory (measuring incompleteness and randomness in mathematics).
1966: Paul Cohen showed that the axiom of choice is unprovable in ZF even without urelements.
1970: Hilbert's tenth problem is proven unsolvable: there is no recursive solution to decide whether a Diophantine equation (multivariable polynomial equation) has a solution in integers.
1971: Suslin's problem is proven to be independent from ZFC.
Toward resolution of the crisis
Starting in 1935, the Bourbaki group of French mathematicians started publishing a series of books to formalize many areas of mathematics on the new foundation of set theory.
The intuitionistic school did not attract many adherents, and it was not until Bishop's work in 1967 that constructive mathematics was placed on a sounder footing.
One may consider that Hilbert's program has been partially completed, so that the crisis is essentially resolved, satisfying ourselves with lower requirements than Hilbert's original ambitions. His ambitions were expressed in a time when nothing was clear: it was not clear whether mathematics could have a rigorous foundation at all.
There are many possible variants of set theory, which differ in consistency strength, where stronger versions (postulating higher types of infinities) contain formal proofs of the consistency of weaker versions, but none contains a formal proof of its own consistency. Thus the only thing we do not have is a formal proof of consistency of whatever version of set theory we may prefer, such as ZF.
In practice, most mathematicians either do not work from axiomatic systems, or if they do, do not doubt the consistency of ZFC, generally their preferred axiomatic system. In most of mathematics as it is practiced, the incompleteness and paradoxes of the underlying formal theories never played a role anyway, and in those branches in which they do or whose formalization attempts would run the risk of forming inconsistent theories (such as logic and category theory), they may be treated carefully.
The development of category theory in the middle of the 20th century showed the usefulness of set theories guaranteeing the existence of larger classes than does ZFC, such as Von Neumann–Bernays–Gödel set theory or Tarski–Grothendieck set theory, albeit that in very many cases the use of large cardinal axioms or Grothendieck universes is formally eliminable.
One goal of the reverse mathematics program is to identify whether there are areas of "core mathematics" in which foundational issues may again provoke a crisis.
| Mathematics | Foundations | null |
169493 | https://en.wikipedia.org/wiki/Petal | Petal | Petals are modified leaves that form an inner whorl surrounding the reproductive parts of flowers. They are often brightly coloured or unusually shaped to attract pollinators. All of the petals of a flower are collectively known as the corolla. Petals are usually surrounded by an outer whorl of modified leaves called sepals, that collectively form the calyx and lie just beneath the corolla. The calyx and the corolla together make up the perianth, the non-reproductive portion of a flower. When the petals and sepals of a flower are difficult to distinguish, they are collectively called tepals. Examples of plants in which the term tepal is appropriate include genera such as Aloe and Tulipa. Conversely, genera such as Rosa and Phaseolus have well-distinguished sepals and petals. When the undifferentiated tepals resemble petals, they are referred to as "petaloid", as in petaloid monocots, orders of monocots with brightly coloured tepals. Since they include Liliales, an alternative name is lilioid monocots.
Although petals are usually the most conspicuous parts of animal-pollinated flowers, wind-pollinated species, such as the grasses, either have very small petals or lack them entirely (apetalous).
Corolla
The collection of all petals in a flower is referred to as the corolla. The role of the corolla in plant evolution has been studied extensively since Charles Darwin postulated a theory of the origin of elongated corollae and corolla tubes.
A corolla of separate petals, without fusion of individual segments, is apopetalous. If the petals are free from one another in the corolla, the plant is polypetalous or choripetalous; while if the petals are at least partially fused, it is gamopetalous or sympetalous. In the case of fused tepals, the term is syntepalous. The corolla in some plants forms a tube.
Variations
Petals can differ dramatically in different species. The number of petals in a flower may hold clues to a plant's classification. For example, flowers on eudicots (the largest group of dicots) most frequently have four or five petals while flowers on monocots have three or six petals, although there are many exceptions to this rule.
The petal whorl or corolla may be either radially or bilaterally symmetrical. If all of the petals are essentially identical in size and shape, the flower is said to be regular or actinomorphic (meaning "ray-formed"). Many flowers are symmetrical in only one plane (i.e., symmetry is bilateral) and are termed irregular or zygomorphic (meaning "yoke-" or "pair-formed"). In irregular flowers, other floral parts may be modified from the regular form, but the petals show the greatest deviation from radial symmetry. Examples of zygomorphic flowers may be seen in orchids and members of the pea family.
In many plants of the aster family such as the sunflower, Helianthus annuus, the circumference of the flower head is composed of ray florets. Each ray floret is anatomically an individual flower with a single large petal. Florets in the centre of the disc typically have no or very reduced petals. In some plants such as Narcissus, the lower part of the petals or tepals are fused to form a floral cup (hypanthium) above the ovary, and from which the petals proper extend.
A petal often consists of two parts: the upper broader part, similar to a leaf blade, also called the blade; and the lower narrower part, similar to a leaf petiole, called the claw, separated from each other at the limb. Claws are distinctly developed in petals of some flowers of the family Brassicaceae, such as Erysimum cheiri.
The inception and further development of petals show a great variety of patterns. Petals of different species of plants vary greatly in colour or colour pattern, both in visible light and in ultraviolet. Such patterns often function as guides to pollinators and are variously known as nectar guides, pollen guides, and floral guides.
Genetics
The genetics behind the formation of petals, in accordance with the ABC model of flower development, are that sepals, petals, stamens, and carpels are modified versions of each other. It appears that the mechanisms to form petals evolved very few times (perhaps only once), rather than evolving repeatedly from stamens.
Significance of pollination
Pollination is an important step in the sexual reproduction of higher plants. Pollen is produced by the male flower or by the male organs of hermaphroditic flowers.
Pollen does not move on its own and thus requires wind or animal pollinators to disperse the pollen to the stigma of the same or nearby flowers. However, pollinators are rather selective in determining the flowers they choose to pollinate. This develops competition between flowers and as a result flowers must provide incentives to appeal to pollinators (unless the flower self-pollinates or is involved in wind pollination). Petals play a major role in competing to attract pollinators. Henceforth pollination dispersal could occur and the survival of many species of flowers could prolong.
Functions and purposes
Petals have various functions and purposes depending on the type of plant. In general, petals operate to protect some parts of the flower and attract/repel specific pollinators.
Function
This is where the positioning of the flower petals are located on the flower is the corolla e.g. the buttercup having shiny yellow flower petals which contain guidelines amongst the petals in aiding the pollinator towards the nectar. Pollinators have the ability to determine specific flowers they wish to pollinate. Using incentives, flowers draw pollinators and set up a mutual relation between each other in which case the pollinators will remember to always guard and pollinate these flowers (unless incentives are not consistently met and competition prevails).
Scent
The petals could produce different scents to allure desirable pollinators or repel undesirable pollinators. Some flowers will also mimic the scents produced by materials such as decaying meat, to attract pollinators to them.
Colour
Various colour traits are used by different petals that could attract pollinators that have poor smelling abilities, or that only come out at certain parts of the day. Some flowers can change the colour of their petals as a signal to mutual pollinators to approach or keep away.
Shape and size
Furthermore, the shape and size of the flower/petals are important in selecting the type of pollinators they need. For example, large petals and flowers will attract pollinators at a large distance or that are large themselves.
Collectively, the scent, colour, and shape of petals all play a role in attracting/repelling specific pollinators and providing suitable conditions for pollinating. Some pollinators include insects, birds, bats, and wind.
In some petals, a distinction can be made between a lower narrowed, stalk-like basal part referred to as the claw, and a wider distal part referred to as the blade (or limb). Often, the claw and blade are at an angle with one another.
Types of pollination
Wind pollination
Wind-pollinated flowers often have small, dull petals and produce little or no scent. Some of these flowers will often have no petals at all. Flowers that depend on wind pollination will produce large amounts of pollen because most of the pollen scattered by the wind tends to not reach other flowers.
Attracting insects
Flowers have various regulatory mechanisms to attract insects. One such helpful mechanism is the use of colour guiding marks. Insects such as the bee or butterfly can see the ultraviolet marks which are contained on these flowers, acting as an attractive mechanism which is not visible towards the human eye. Many flowers contain a variety of shapes acting to aid with the landing of the visiting insect and also influence the insect to brush against anthers and stigmas (parts of the flower). One such example of a flower is the pohutukawa (Metrosideros excelsa), which acts to regulate colour in a different way. The pohutukawa contains small petals also having bright large red clusters of stamens.
Another attractive mechanism for flowers is the use of scents which are highly attractive to humans. One such example is the rose. On the other hand, some flowers produce the smell of rotting meat and are attractive to insects such as flies. Darkness is another factor that flowers have adapted to as nighttime conditions limit vision and colour-perception. Fragrancy can be especially useful for flowers that are pollinated at night by moths and other flying insects.
Attracting birds
Flowers are also pollinated by birds and must be large and colourful to be visible against natural scenery. In New Zealand, such bird–pollinated native plants include: kowhai (Sophora species), flax (Phormium tenax) and kaka beak (Clianthus puniceus). Flowers adapt the mechanism on their petals to change colour in acting as a communicative mechanism for the bird to visit. An example is the tree fuchsia (Fuchsia excorticata), which are green when needing to be pollinated and turn red for the birds to stop coming and pollinating the flower.
Bat-pollinated flowers
Flowers can be pollinated by short-tailed bats. An example of this is the dactylanthus (Dactylanthus taylorii). This plant has its home under the ground acting the role of a parasite on the roots of forest trees. The dactylanthus has only its flowers pointing to the surface and the flowers lack colour but have the advantage of containing much nectar and a strong scent. These act as a useful mechanism in attracting the bat.
| Biology and health sciences | Plant anatomy and morphology: General | Biology |
169498 | https://en.wikipedia.org/wiki/Alopecia%20areata | Alopecia areata | Alopecia areata, also known as spot baldness, is a condition in which hair is lost from some or all areas of the body. It often results in a few bald spots on the scalp, each about the size of a coin. Psychological stress and illness are possible factors in bringing on alopecia areata in individuals at risk, but in most cases there is no obvious trigger. People are generally otherwise healthy. In a few cases, all the hair on the scalp is lost (alopecia totalis), or all body hair is lost (alopecia universalis). Hair loss can be permanent, or temporary.
Alopecia areata is believed to be an autoimmune disease resulting from a breach in the immune privilege of the hair follicles. Risk factors include a family history of the condition. Among identical twins, if one is affected, the other has about a 50% chance of also being affected. The underlying mechanism involves failure by the body to recognize its own cells, with subsequent immune-mediated destruction of the hair follicle.
No cure for the condition is known. Some treatments, particularly triamcinolone injections and 5% minoxidil topical creams, are effective in speeding hair regrowth. Sunscreen, head coverings to protect from cold and sun, and glasses, if the eyelashes are missing, are also recommended. In more than 50% of cases of sudden-onset localized "patchy" disease, hair regrows within a year. In patients with only one or two patches, this one-year recovery will occur in up to 80%. However, many people will have more than one episode over the course of a lifetime. In many patients, hair loss and regrowth occurs simultaneously over the course of several years. Among those in whom all body hair is lost, fewer than 10% recover.
About 0.15% of people are affected at any one time, and 2% of people are affected at some point in time. Onset is usually in childhood. Females are affected at higher rates than males.
Signs and symptoms
Typical first symptoms of alopecia areata are small bald patches. The underlying skin is unscarred and looks superficially normal. Although these patches can take many shapes, they are usually round or oval. Alopecia areata most often affects the scalp and beard, but may occur on any part of the body with hair. Different areas of the skin may exhibit hair loss and regrowth at the same time. The disease may also go into remission for a time, or may be permanent. It is common in children.
The area of hair loss may tingle or be mildly painful. The hair tends to fall out over a short period of time, with the loss commonly occurring more on one side of the scalp than the other.
Exclamation point hairs, narrower along the length of the strand closer to the base, producing a characteristic "exclamation point" appearance, are often present. These hairs are very short (3–4 mm), and can be seen surrounding the bald patches.
When healthy hair is pulled out, at most a few should come out, and ripped hair should not be distributed evenly across the tugged portion of the scalp. In cases of alopecia areata, hair tends to pull out more easily along the edge of the patch where the follicles are already being attacked by the body's immune system than away from the patch where they are still healthy.
Nails may have pitting or trachyonychia. Onychoptosis defluvium, also known as alopecia unguium, is casting off the nail seen in association with alopecia areata.
Causes
Alopecia areata is thought to be a systemic autoimmune disorder in which the body attacks its own anagen hair follicles and suppresses or stops hair growth. For example, T cell lymphocytes cluster around affected follicles, causing inflammation and subsequent hair loss. Hair follicles in a normal state are thought to be kept secure from the immune system, a phenomenon called immune privilege. A breach in this immune privilege state is considered as the cause of alopecia areata. A few cases of babies being born with congenital alopecia areata have been reported. It is recognized as a type 1 inflammatory disease.
Alopecia areata is not contagious. It occurs more frequently in people who have affected family members, suggesting heredity may be a factor. Strong evidence of genetic association with increased risk for alopecia areata was found by studying families with two or more affected members. This study identified at least four regions in the genome that are likely to contain these genes. In addition, alopecia areata shares genetic risk factors with other autoimmune diseases, including rheumatoid arthritis, type 1 diabetes, and celiac disease. It may be the only manifestation of celiac disease.
Endogenous retinoids metabolic defect is a key part of the pathogenesis of the alopecia areata.
In 2010, a genome-wide association study was completed that identified 129 single nucleotide polymorphisms that were associated with alopecia areata. The genes that were identified include those involved in
controlling the activation and proliferation of regulatory T cells, cytotoxic T lymphocyte-associated antigen 4, interleukin-2, interleukin-2 receptor A, and Eos (also known as Ikaros family zinc finger 4), as well as the human leukocyte antigen. The study also identified two genes, PRDX5 and STX17, that are expressed in the hair follicle.
A psychodermatological connection is noted with impairment in psychiatric comorbidities including mental well-being, self esteem and mental disorders acting as pathogenic triggers for alopecia areata.
Diagnosis
Alopecia areata is usually diagnosed based on clinical features.
Trichoscopy may aid in establishing the diagnosis. In alopecia areata, trichoscopy shows regularly distributed "yellow dots" (hyperkeratotic plugs), small exclamation-mark hairs, and "black dots" (destroyed hairs in the hair follicle opening).
Oftentimes, however, discrete areas of hair loss surrounded by exclamation mark hairs is sufficient for clinical diagnosis of alopecia areata. Sometimes, reddening of the skin, erythema, may also be present in the balding area.
A biopsy is rarely needed to make the diagnosis or aid in the management of alopecia areata. Histologic findings may include peribulbar lymphocytic infiltration resembling a "swarm of bees", a shift in the anagen-to-telogen ratio towards telogen, and dilated follicular infundibulae. Other helpful findings can include pigment incontinence in the hair bulb and follicular stelae. Occasionally, in inactive alopecia areata, no inflammatory infiltrates are found.
Classification
Commonly, alopecia areata involves hair loss in one or more round spots on the scalp.
Hair may also be lost more diffusely over the whole scalp, in which case the condition is called diffuse alopecia areata.
Alopecia areata monolocularis describes baldness in only one spot. It may occur anywhere on the head.
Alopecia areata multilocularis refers to multiple areas of hair loss.
Ophiasis refers to hair loss in the shape of a wave at the circumference of the head.
The disease may be limited only to the beard, in which case it is called alopecia areata barbae.
If the person loses all the hair on the scalp, the disease is then called alopecia areata totalis.
If all body hair, including pubic hair, is lost, the diagnosis then becomes alopecia areata universalis.
Alopecia areata totalis and universalis are rare.
Treatment
The objective assessment of treatment efficacy is very difficult and spontaneous remission is unpredictable, but if the affected area is patchy, the hair may regrow spontaneously in many cases. None of the existing therapeutic options are curative or preventive. A 2020 systematic review showed greater than 50% hair regrowth in 80.9% of patients treated with 5 mg/mL triamcinolone injections. A Cochrane-style systematic review published in 2019 showed 5% topical minoxidil was more than 8x more associated with >50% hair regrowth at 6 months compared to placebo. In cases of severe hair loss, limited success has been achieved by using the corticosteroid medications clobetasol or fluocinonide as an injection or cream. Application of corticosteroid creams to the affected skin is less effective and takes longer to produce results. Steroid injections are commonly used in sites where the areas of hair loss on the head are small or especially where eyebrow hair has been lost. Whether they are effective is uncertain. Some other medications that have been used are minoxidil, Elocon (mometasone) ointment (steroid cream), irritants (anthralin or topical coal tar), and topical immunotherapy ciclosporin, sometimes in different combinations. Topical corticosteroids frequently fail to enter the skin deeply enough to affect the hair bulbs, which are the treatment target, and small lesions typically also regrow spontaneously. Oral corticosteroids may decrease the hair loss, but only for the period during which they are taken, and these medications can cause serious side effects. No one treatment is effective in all cases, and some individuals may show no response to any treatment.
For more severe cases, studies have shown promising results with the individual use of the immunosuppressant methotrexate or adjunct use with corticosteroids.
When alopecia areata is associated with celiac disease, treatment with a gluten-free diet allows for complete and permanent regrowth of scalp and other body hair in many people, but in others, remissions and recurrences are seen. This improvement is probably due to the normalization of the immune response as a result of gluten withdrawal from the diet.
In June 2022, the U.S. Food and Drug Administration (FDA) authorized baricitinib, a Janus kinase (JAK) inhibitor, for the treatment of severe alopecia areata.
Ritlecitinib (Litfulo) was approved for medical use in the United States in June 2023.
Fecal matter transplants (FMT) have been shown to reverse AA and support hair growth, with long lasting results, even going as far as growing additional hair on arms and face while grey hairs even regained colour. This supports the idea of a connection between gut microbiota having a part in hair loss.
Hair transplantation may be an alternative for patients with chronic local alopecia areata. The fact that the disease is autoimmune and progresses with relapses is one of the biggest question marks before surgery. There have been case reports in the literature since the early 2000s. However, in an article published long-term follow-up; It is reported that the hair transplanted to the eyebrow area falls out again due to the recurrence of the disease. A similar situation was not mentioned in previous studies on this subject. Perhaps the long-term follow-ups of other studies were not sufficient.
Deuruxolitinib (Leqselvi) was approved for medical use in the United States in July 2024.
Prognosis
In most cases that begin with a small number of patches of hair loss, hair grows back after a few months to a year. In cases with a greater number of patches, hair can either grow back or progress to alopecia areata totalis or, in rare cases, alopecia areata universalis.
No loss of body function occurs, and the effects of alopecia areata are psychological (loss of self-image due to hair loss), although these can be severe. Loss of hair also means the scalp burns more easily in the sun. Patients may also have aberrant nail formation because keratin forms both hair and nails.
Hair may grow back and then fall out again later. This may not indicate a recurrence of the condition, but rather a natural cycle of growth-and-shedding from a relatively synchronised start; such a pattern will fade over time. Episodes of alopecia areata before puberty predispose to chronic recurrence of the condition.
Alopecia can be the cause of psychological stress. Because hair loss can lead to significant changes in appearance, individuals with it may experience social phobia, anxiety, and depression.
Epidemiology
The condition affects 0.1%–0.2% of the population, with a lifetime risk of 1%-2%, and is more common in females.
Alopecia areata occurs in people who are otherwise healthy and have no other skin disorders. Initial presentation most commonly occurs in the early childhood, late teenage years, or young adulthood, but can happen at any ages. Patients also tend to have a slightly higher incidence of conditions related to the immune system, such as asthma, allergies, atopic dermatitis, and hypothyroidism.
Society and culture
The term alopecia, used by physicians dating back to Hippocrates, originates from the Greek word for fox, "alopex", and was so-named due to fur loss seen in fox mange. "Areata" is derived from the Latin word, "area", meaning a vacant space or patch.
Alopecia areata and alopecia barbae have been identified by some as the biblical condition that is part of the greater family of skin disorders; the said disorders are purported to being discussed in the Book of Leviticus, chapter 13.
Notable people
NASCAR driver Joey Logano, obstacle athlete Kevin Bull, politicians Peter Dutton and Ayanna Pressley, K-pop singer Peniel Shin of BtoB, actors Christopher Reeve, Anthony Carrigan and Alan Fletcher, and actresses Jada Pinkett Smith, May Calamawy, and Lili Reinhart all have some form of alopecia areata.
Research
Many medications are being studied.
In 2014, preliminary findings showing that oral ruxolitinib, a drug approved by the US Food and Drug Administration (FDA) for bone marrow disorder myelofibrosis, restored hair growth in three individuals with long-standing and severe disease.
In March 2020, the US FDA granted breakthrough therapy designation to baricitinib for the systematic treatment of alopecia areata and granted approval in June 2022, with a 32% efficacy rate for people with 50% hair loss reaching 80% scalp coverage in 36 weeks. It acts as an inhibitor of janus kinase (JAK), blocking the subtypes JAK1 and JAK2.
| Biology and health sciences | Health and fitness: General | Health |
169546 | https://en.wikipedia.org/wiki/Suede | Suede | Suede (pronounced ) is a type of leather with a fuzzy, napped finish, commonly used for jackets, shoes, fabrics, purses, furniture, and other items.
Suede is made from the underside of the animal skin, which is softer and more pliable than the outer skin layer, though not as durable.
Etymology
The term comes from the French , which literally means "gloves from Sweden". The term was first used by The Oxford English Dictionary in 1884.
Production
Suede leather is made from the underside of the skin, primarily from lamb, although goat, calf, and deer are commonly used. Splits from thick hides of cow and deer are also sueded, but, due to the fiber content, have a shaggy nap.
Characteristics
Because suede does not include the tough exterior skin layer, it is less durable, but softer, than the standard "full-grain" leather. Its softness, thinness, and pliability make it suitable for clothing and delicate uses; suede was originally used for women's gloves, hence its etymology (see above). Suede leather is also popular in upholstery, shoes, bags, and other accessories, and as a lining for other leather products. Due to its textured nature and open pores, suede may become dirty and quickly absorb liquids. Suede is often used in place of leather when more breathability (air permeation) is needed such as with hot weather footgear.
Preservation and conditioning
A variety of environmental factors including salt, dirt, water, oils and moisture can stain or wear out suede. Since excess moisture can damage suede, it should not be cleaned with soap and water or machine washed. Suede brushes and suede rubbers, as well as a nail files, are tools that may be used to clean suede, often in conjunction with white vinegar or cornstarch. Suede protector spray can be applied after cleaning to preserve the integrity of the fabric longer.
In popular culture
Suede's absorbent nature was highlighted in the Seinfeld episode "The Jacket", in which Jerry ventures outside into the snow and ruins his exorbitantly priced suede jacket.
"Blue Suede Shoes" is a well-known early rock-n-roll song written by Carl Perkins and also covered by Elvis Presley.
"Weird Al" Yankovic wrote and performed the song "King of Suede".
"Suedehead" A skinhead subculture and song by English singer/songwriter Morrissey.
" Johnny Suede" A film starring Brad Pitt where he plays a down and out musician with a huge pompadour haircut.
| Technology | Materials | null |
169552 | https://en.wikipedia.org/wiki/Fifth%20force | Fifth force | In physics, a fifth force refers to a hypothetical fundamental interaction (also known as fundamental force) beyond the four known interactions in nature: gravitational, electromagnetic, strong nuclear, and weak nuclear forces.
Some speculative theories have proposed a fifth force to explain various anomalous observations that do not fit existing theories. The specific characteristics of a putative fifth force depend on which hypothesis is being advanced. No evidence to support these models has been found.
The term is also used as "the Fifth force" when referring to a specific theory advanced by Ephraim Fischbach in 1971 to explain experimental deviations in the theory of gravity. Later analysis failed to reproduce those deviations.
History
The term fifth force originates in a 1986 paper by Ephraim Fischbach et al. who reanalyzed the data from the Eötvös experiment of Loránd Eötvös from earlier in the century; the reanalysis found a distance dependence to gravity that deviates from the inverse square law.
The reanalysis was sparked by theoretical work in 1971 by Fujii proposing a model that changes distance dependence with a Yukawa potential-like term:
The parameter characterizes the strength and the range of the interaction. Fischbach's paper found a strength around 1% of gravity and a range of a few hundred meters.
The effect of this potential can be described equivalently as exchange of vector and/or scalar bosons, that is a predicting as yet undetected new particles.
However, many subsequent attempts to reproduce the deviations have failed.
Theory
Theoretical proposals for a fifth-force are driven by inconsistencies between the existing models of general relativity and quantum field theory, and also between the hierarchy problem and the cosmological constant problem. Both issues suggest the possibility of corrections to the gravitational potential around .
The accelerating expansion of the universe has been attributed to a form of energy called dark energy. Some physicists speculate that a form of dark energy called quintessence could be a fifth force.
Experimental approaches
There are at least three kinds of searches that can be undertaken, which depend on the kind of force being considered, and its range.
Equivalence principle
One way to search for a fifth force is with tests of the strong equivalence principle, one of the most powerful tests of general relativity, also known as Einstein's theory of gravity. Alternative theories of gravity, such as Brans–Dicke theory, postulate a fifth possibly one with infinite range. This is because gravitational interactions, in theories other than general relativity, have degrees of freedom other than the "metric", which dictates the curvature of space, and different kinds of degrees of freedom produce different effects. For example, a scalar field cannot produce the bending of light rays.
The fifth force would manifest itself in an effect on solar system orbits, called the Nordtvedt effect. This is tested with Lunar Laser Ranging experiment and very-long-baseline interferometry.
Extra dimensions
Another kind of fifth force, which arises in Kaluza–Klein theory, where the universe has extra dimensions, or in supergravity or string theory is the Yukawa force, which is transmitted by a light scalar field (i.e. a scalar field with a long Compton wavelength, which determines the range). This has prompted a much recent interest, as a theory of supersymmetric large extra dimensions with size slightly less than a has prompted an experimental effort to test gravity on very small scales. This requires extremely sensitive experiments which search for a deviation from the inverse-square law of gravity over a range of distances. Essentially, they are looking for signs that the Yukawa interaction is engaging at a certain length.
Australian researchers, attempting to measure the gravitational constant deep in a mine shaft, found a discrepancy between the predicted and measured value, with the measured value being two percent too small. They concluded that the results may be explained by a repulsive fifth force with a range from a few centimetres to a kilometre. Similar experiments have been carried out on board a submarine, USS Dolphin (AGSS-555), while deeply submerged. A further experiment measuring the gravitational constant in a deep borehole in the Greenland ice sheet found discrepancies of a few percent, but it was not possible to eliminate a geological source for the observed signal.
Earth's mantle
Another experiment uses the Earth's mantle as a giant particle detector, focusing on geoelectrons.
Cepheid variables
Jain et al. (2012) examined existing data on the rate of pulsation of over a thousand cepheid variable stars in 25 galaxies. Theory suggests that the rate of cepheid pulsation in galaxies screened from a hypothetical fifth force by neighbouring clusters, would follow a different pattern from cepheids that are not screened. They were unable to find any variation from Einstein's theory of gravity.
Other approaches
Some experiments used a lake plus a tower that is eters high. A comprehensive review by Ephraim Fischbach and Carrick Talmadge suggested there is no compelling evidence for the fifth force, though scientists still search for it. The Fischbach–Talmadge article was written in 1992, and since then, other evidence has come to light that may indicate a fifth force.
The above experiments search for a fifth force that is, like gravity, independent of the composition of an object, so all objects experience the force in proportion to their masses. Forces that depend on the composition of an object can be very sensitively tested by torsion balance experiments of a type invented by Loránd Eötvös. Such forces may depend, for example, on the ratio of protons to neutrons in an atomic nucleus, nuclear spin, or the relative amount of different kinds of binding energy in a nucleus (see the semi-empirical mass formula). Searches have been done from very short ranges, to municipal scales, to the scale of the Earth, the Sun, and dark matter at the center of the galaxy.
Claims of new particles
In 2015, Attila Krasznahorkay at ATOMKI, the Hungarian Academy of Sciences's Institute for Nuclear Research in Debrecen, Hungary, and his colleagues posited the existence of a new, light boson only 34 times heavier than the electron (17 MeV). In an effort to find a dark photon, the Hungarian team fired protons at thin targets of lithium-7, which created unstable beryllium-8 nuclei that then decayed and ejected pairs of electrons and positrons. Excess decays were observed at an opening angle of 140° between the and , and a combined energy of 17 MeV, which indicated that a small fraction of beryllium-8 will shed excess energy in the form of a new particle.
In November 2019, Krasznahorkay announced that he and his team at ATOMKI had successfully observed the same anomalies in the decay of stable helium atoms as had been observed in beryllium-8, strengthening the case for the X17 particle's existence.
Feng et al. (2016) proposed that a protophobic (i.e. "proton-ignoring") X-boson with a mass of 16.7 MeV with suppressed couplings to protons relative to neutrons and electrons and femtometer range could explain the data. The force may explain the muon anomaly and provide a dark matter candidate. Several research experiments are underway to attempt to validate or refute these results.
| Physical sciences | Physics basics: General | Physics |
169553 | https://en.wikipedia.org/wiki/Online%20dating | Online dating | Online dating, also known as internet dating, virtual dating, or mobile app dating, is a method used by people with a goal of searching for and interacting with potential romantic or sexual partners, via the internet. An online dating service is a company that promotes and provides specific mechanisms for the practice of online dating, generally in the form of dedicated websites or software applications accessible on personal computers or mobile devices connected to the internet. A wide variety of unmoderated matchmaking services, most of which are profile-based with various communication functionalities, is offered by such companies.
Online dating services allow users to become "members" by creating a profile and uploading personal information including (but not limited to) age, gender, sexual orientation, location, and appearance. Most services also encourage members to add photos or videos to their profile. Once a profile has been created, members can view the profiles of other members of the service, using the visible profile information to decide whether or not to initiate contact. Most services offer digital messaging, while others provide additional services such as webcasts, online chat, telephone chat (VOIP), and message boards. Members can constrain their interactions to the online space, or they can arrange a date to meet in person.
A great diversity of online dating services currently exist. Some have a broad membership base of diverse users looking for many different types of relationships. Other sites target highly specific demographics based on features like shared interests, location, religion, sexual orientation or relationship type. Online dating services also differ widely in their revenue streams. Some sites are completely free and depend on advertising for revenue. Others utilize the freemium revenue model, offering free registration and use, with optional, paid, premium services. Still others rely solely on paid membership subscriptions.
Trends
Social trends and public opinions
A 2005 study found that online daters may have more liberal social attitudes compared to the general population in the United States.
Race and online dating
A 2009 study found that African Americans were the least desired demographic in online dating; and were the least interested in forming interracial relationships with non-Black Americans.
In 2008, a study investigated racial preferences using a sample of 6,070 profiles on Yahoo! Personals. Just 29% of white men excluded women of color, compared to the 64% of white women who excluded men of color. Follow-up studies conducted by these authors in 2009 and 2011 found similar patterns: white women were less open to interracial relationships than white men.
In 2018, a study analyzed the activity of approximately 200,000 users of an online dating app in the United States. The authors found that White men and Asian women were the most desired.
In 2021, a comprehensive analysis of online dating trends in the United States suggested that the rise of online dating has exacerbated underlying racial biases in dating. The authors found that White men were preferred by women of color, while men of color generally preferred women of color. White men were accepting of Asian and Hispanic women, yet White women tended to exclude non-White men.
However, these authors also disputed some common notions about racial bias in online dating. For example, White women did not reject Asian men more so than Black or Hispanic men. Black and Hispanic women were just as accepting of Asian men as they were of men of the same race. This is inconsistent with the idea that Asian men are particularly disadvantaged in online dating, relative to other men of color. The authors also dispute the notion that Asian women's high outmarriage rate is due to "self hatred", as their interviews found that these marriages form out of perceived compatibility, rather than self hatred. Gay Hispanic men did not have a preference for white partners.
Gender
According to a 2018 study, among American daters, male desirability increased until the age of 50; while women's desirability declined steeply after the age of 20.
In terms of educational attainment, men's desirability only increased the more educated they were. For women, however, educational attainment beyond the level of a bachelor's degree actually decreased their desirability. The authors suggested that besides individual preferences and partner availability, this pattern may be due to the fact that by the late 2010s, women were more likely to attend and graduate from university. In order to estimate the desirability of a given individual, the researchers looked at the number of messages they received and the desirability of the senders.
Developmental psychologist Michelle Drouin, who was not involved in the study, told The New York Times this finding is in accordance with theories in psychology and sociology based on biological evolution in that youth is a sign of fertility. She added that women with advanced degrees are often viewed as more focused on their careers than family. Licensed psychotherapist Stacy Kaiser told MarketWatch men typically prefer younger women because "they are more easy to impress; they are more (moldable) in terms of everything from emotional behavior to what type of restaurant to eat at," and because they tend to be "more fit, have less expectations and less baggage." On the other hand, women look for (financial) stability and education, attributes that come with age, said Kaiser. These findings regarding age and attractiveness are consistent with earlier research by the online dating services OKCupid and Zoosk.
In 2016, Gareth Tyson of the Queen Mary University of London and his colleagues published a paper analyzing the behavior of Tinder users in New York City and London. In order to minimize the number of variables, they created profiles of white heterosexual people only. For each sex, there were three accounts using stock photographs, two with actual photographs of volunteers, one with no photos whatsoever, and one that was apparently deactivated. The researchers pointedly only used pictures of people of average physical attractiveness. Tyson and his team wrote an algorithm that collected the biographical information of all the matches, liked them all, then counted the number of returning likes.
They found that men and women employed drastically different mating strategies. Men liked a large proportion of the profiles they viewed, but received returning likes only 0.6% of the time; women were much more selective but received matches 10% of the time. Men received matches at a much slower rate than women. Once they received a match, women were far more likely than men to send a message, 21% compared to 7%, but they took more time before doing so. Tyson and his team found that for the first two-thirds of messages from each sex, women sent them within 18 minutes of receiving a match compared to five minutes for men. Men's first messages had an average of a dozen characters, and were typical simple greetings; by contrast, initial messages by women averaged 122 characters.
Tyson and his collaborators found that the male profiles that had three profile pictures received far more matches than those without one. By sending out questionnaires to frequent Tinder users, the researchers discovered that the reason why men tended to like a large proportion of the women they saw was to increase their chances of getting a match. This led to a feedback loop in which men liked more and more of the profiles they saw while women could afford to be even more selective in liking profiles because of a greater probability of a match.
Aided by the text-analysis program Linguistic Inquiry and Word Count, Bruch and Newman discovered that men generally had lower chances of receiving a response after sending more "positively worded" messages. When a man tried to woo a woman more desirable than he was, he received a response 21% of the time; by contrast, when a woman attempted to court a man, she received a reply about half the time. In fact, over 80% of the first messages in the data set obtained for the purposes of the study were from men, and women were highly selective in choosing whom to respond to, a rate of less than 20%. Therefore, studying women's replies yielded much insight into their preferences. Bruch and Newman were also able to establish the existence of dating 'leagues'. Generally speaking, people were able to accurately estimate where they ranked on the dating hierarchy. Very few responded to the messages of people less desirable than they were. Nevertheless, although the probability of a response is low, it is well above zero, and if the other person does respond, it can a self-esteem booster, said Kaiser. Co-author of the study Mark Newman told BBC News, "There is a trade-off between how far up the ladder you want to reach and how low a reply rate you are willing to put up with." Bruch and Newman found that while people spent a lot of time crafting lengthy messages to those they considered to be a highly desirable partner, this hardly made a difference, judging by the response rate. Keeping messages concise is well-advised. Previous studies also suggest that about 70% of the dating profile should be about oneself and the rest about the desired partner.
Data from the Chinese online dating giant Zhenai.com reveals that while men are most interested in how a woman looks, women care more about a man's income. Profession is also quite important. Chinese men favor women working as primary school teachers and nurses while Chinese women prefer men in the IT or finance industry. Women in IT or finance are the least desired. Zhenai enables users to send each other digital "winks". For a man, the more money he earns the more "winks" he receives. For a woman, her income does not matter until the 50,000-yuan mark (US$7,135), after which the number of "winks" falls slightly. Men typically prefer women three years younger than they are whereas women look for men who are three years older on average. However, this changes if the man becomes exceptionally wealthy; the more money he makes the more likely he is to look for younger women.
In general, people in their 20s employ the "self-service dating service" while women in their late 20s and up tend to use the matchmaking service. This is because of the social pressure in China on "leftover women" (Sheng nu), meaning those in their late 20s but still not married. Women who prefer not to ask potentially embarrassing questions – such as whether both spouses will handle household finances, whether or not they will live with his parents, or how many children he wants to have, if any – will get a matchmaker to do it for them. Both sexes prefer matchmakers who are women.
Desirability and physical appearance
At least three quarters of the sample surveyed attempted to date aspirationally, meaning they tried to initiate a relationship with someone who was more desirable, 25% more desirable, to be exact. Bruch recommended sending out more greeting messages, noting that people sometimes managed to upgrade their 'league'. Michael Rosenfeld, a sociologist not involved with the study, told The Atlantic, "The idea that persistence pays off makes sense to me, as the online-dating world has a wider choice set of potential mates to choose from. The greater choice set pays dividends to people who are willing to be persistent in trying to find a mate." Using optimal stopping theory, one can show that the best way to select the best potential partner is to reject the first 37%, then pick the one who is better than the previous set. The probability of picking the best potential mate this way is 37%. (This is approximately the reciprocal of Euler's number, . See derivation of the optimal policy.) However, making online contact is only the first step, and indeed, most conversations failed to birth a relationship. As two potential partners interact more and more, the superficial information available from a dating website or smartphone application becomes less important than their characters.
Despite being a platform designed to be less centered on physical appearance, OkCupid co-founder Christian Rudder stated in 2009 that the male OkCupid users who were rated most physically attractive by female OkCupid users received 11 times as many messages as the lowest-rated male users did, the medium-rated male users received about four times as many messages, and the one-third of female users who were rated most physically attractive by the male users received about two-thirds of all messages sent by male users. According to a former company product manager, the majority of female Bumble users typically set a floor height of six feet for male users which limits their matching opportunities to only 15% of the male population.
Niche dating sites
Sites with specific demographics have become popular as a way to narrow the pool of potential matches. Successful niche sites pair people by race, sexual orientation or religion. In March 2008, the top 5 overall sites held 7% less market share than they did one year ago while the top sites from the top five major niche dating categories made considerable gains. Niche sites cater to people with special interests, such as sports fans, racing and automotive fans, medical or other professionals, people with political or religious preferences, people with medical conditions, or those living in rural farm communities.
Some dating services have been created specifically for those living with HIV and other venereal diseases in an effort to eliminate the need to lie about one's health in order to find a partner. Public health officials in Rhode Island and Utah claimed in 2015 that Tinder and similar apps were responsible for uptick of such conditions.
Some sites, referred to as adult dating sites, match individuals seeking short-term sexual encounters.
Economic trends
Although some sites offer free trials and/or profiles, most memberships can cost upwards of $60 per month. In 2008, online dating services in the United States generated $957 million in revenue.
Most free dating websites depend on advertising revenue, using tools such as Google AdSense and affiliate marketing. Since advertising revenues are modest compared to membership fees, this model requires numerous page views to achieve profitability. However, Sam Yagan describes dating sites as ideal advertising platforms because of the wealth of demographic data made available by users.
In November 2023, the stock prices of Match Group and Bumble were down 31% and 35% on the year respectively, continuing a more than two-year decline since the latter's initial public offering in February 2021 and after posting declines more than double that of the S&P 500 during the 2022 stock market decline. In addition to price increases, slowing paid user growth, and flattening app download rates following the end of the COVID-19 lockdowns, assessments among financial analysts of an oversaturated market, concerns about low consumer satisfaction with the services, and growing skepticism about dating app features and algorithms contributed to the declines. Match Group and Bumble account for nearly the entire market share of the online dating industry, and the companies lost a combined $40 billion in market value from 2021 through 2024. Match Group and Bumble shares continued to fall during the first quarter of 2024 while the S&P 500 rose, and the number of paid users for Match Group fell by 6% during the first quarter of 2024 while Bumble's paid users grew by 18% in comparison to a 3% decline and a 31% increase for each company respectively during the first quarter of 2023.
Matching and divorce rates
In 2012, social psychologists Benjamin Karney, Harry Reis, and others published an analysis of online dating in Psychological Science in the Public Interest that concluded that the matching algorithms of online dating services are only negligibly better at matching people than if they were matched at random. In 2014, Kang Zhao at the University of Iowa constructed a new approach based on the algorithms used by Amazon and Netflix, based on recommendations rather than the autobiographical notes of match seekers. Users' activities reflect their tastes and attractiveness, or the lack thereof, they reasoned. This algorithm increases the chances of a response by 40%, the researchers found. E-commerce firms also employ this "collaborative filtering" technique. Nevertheless, it is still not known what the algorithm for finding the perfect match would be.
However, while collaborative filtering and recommender systems have been demonstrated to be more effective than matching systems based on similarity and complementarity, they have also been demonstrated to be highly skewed to the preferences of early users and against racial minorities such as African Americans and Hispanic Americans which led to the rise of niche dating sites for those groups. In 2014, the Better Business Bureau's National Advertising Division criticized eHarmony's claims of creating a greater number of marriages and more durable and satisfying marriages than alternative dating websites, and in 2018, the Advertising Standards Authority banned eHarmony advertisements in the United Kingdom after the company was unable to provide any evidence to verify its advertisements' claims that its website's matching algorithm was scientifically proven to give its users a greater chance of finding long-term intimate relationships.
Data released by Tinder in 2018 showed that of the 1.6 billion swipes it recorded per day, only 26 million result in matches (a match rate of approximately only 1.63%), despite users logging into the app on average 11 times per day, with male user sessions averaging 7.2 minutes and female user sessions averaging 8.5 minutes (or 79.2 minutes and 93.5 minutes per day respectively). Also, a Tinder user interviewed anonymously in an article published in the December 2018 issue of The Atlantic estimated that only one in 10 of their matches actually resulted in an exchange of messages with the other user they were matched with, with another anonymous Tinder user saying, "Getting right-swiped is a good ego boost even if I have no intention of meeting someone."
In 2012, Karney, Reis, and their co-authors suggested that the availability of a large pool of potential partners "may lead online daters to objectify potential partners and might even undermine their willingness to commit to one of them." In October 2019, a Pew Research Center survey of 4,860 U.S. adults showed that 54 percent of U.S. adults believed that relationships formed through dating sites or apps were just as successful as those that began in person, 38 percent believed these relationships were less successful, while only 5 percent believed them to be more successful.
Noting the research of Karney, Reis, and their co-authors comparing online to offline dating and the research of communications studies scholar Nicole Ellison and her co-authors comparing online dating to comparative shopping, political scientist Robert D. Putnam cited the October 2019 Pew Research Center survey in the afterword to the second edition of Bowling Alone (2020) in expressing skepticism about whether online dating was leading to a greater number of long-term intimate relationships. Social psychologist David Buss has estimated that approximately 30 percent of the men on Tinder are married.
Buss has argued further "Apps like Tinder and OkCupid give people the impression that there are thousands or millions of potential mates out there. One dimension of this is the impact it has on men's psychology. When there is ... a perceived surplus of women, the whole mating system tends to shift towards short-term dating," and there is a feeling of disconnect when choosing future partners. In addition, the cognitive process identified by psychologist Barry Schwartz as the "paradox of choice" (also referred to as "choice overload" or "fear of a better option") was cited in an article published in The Atlantic that suggested that the appearance of an abundance of potential partners causes online daters to be less likely to choose a partner and be less satisfied with their choices of partners.
Research on associations between online dating and divorce rates have found conflicting results. While research published in the Journal of Family and Economic Issues in September 2011 found no relationship between increased internet access and higher divorce rates in the United States, subsequent research published in the Review of Economics of the Household in June 2020 did find a correlation between increased access to broadband internet or mobile phones and higher divorce rates in rural counties and lower divorces rates in metropolitan areas in the United States. In June 2013, PNAS USA published a representative survey of 19,131 U.S. adults married between 2005 and 2012 that found that marriages that began online were slightly less likely to result in separation or divorce in comparison to marriages formed offline and were associated with slightly higher marital satisfaction. In July 2014, Computers in Human Behavior published a study that found that after controlling for various economic, demographic, and psychological variables that state-by-state differences in the United States in Facebook and other social networking service (SNS) user account rates was correlated with higher divorce rates and diminished marriage quality. In October 2015, Cyberpsychology, Behavior, and Social Networking published a study of 371 undergraduate students at a university in the Midwestern United States that found that Facebook friend lists increased physical and emotional infidelities among couples, lowered relationship commitment, and diminished relationship quality due to psychological priming effects.
In November 2016, the Journal of International Social Issues published a study that found that U.S. states with a higher Google Trends search volume index for Match.com in 2013 had fewer marriages in 2014, while U.S. states with higher search volume indices for Hinge, Bumble, Plenty of Fish, and Facebook in 2013 had a greater number of divorces in 2014. In February 2019, Technological Forecasting and Social Change published a study examining associations between broadband internet access and divorce in China using provincial data from 2002 to 2014 that found that for every 1% increase in the number of broadband subscribers the number of divorces grew by 0.008%. In December 2020, PLOS One published a study on online dating in Switzerland that found that couples formed through online dating had stronger cohabiting intentions than those formed offline and no differences in relationship satisfaction. In January 2024, Computers in Human Behavior published a survey of 923 married U.S. adults where roughly half of the subjects met their spouses online that found evidence for an "online dating effect" where online daters reported less satisfying and durable marriages, but the researchers suggested that the differences could be explained by societal marginalization and geographic distance.
Online matchmaking services
In 2008, a variation of the online dating model emerged in the form of introduction sites, where members have to search and contact other members, who introduce them to other members whom they deem compatible. Introduction sites differ from the traditional online dating model, and attracted many users and significant investor interest.
In China, the number of separations per a thousand couples doubled, from 1.46 in 2006 to about three in 2016, while the number of actual divorces continues to rise, according to the Ministry of Civil Affairs. Demand for online dating services among divorcees keeps growing, especially in the large cities such as Beijing, Shanghai, Shenzhen and Guangzhou. In addition, more and more people are expected to use online dating and matchmaking services as China continues to urbanize in the late 2010s and 2020s.
Reception
Opinions and usage of online dating services also differ widely. A 2005 study of data collected by the Pew Internet & American Life Project found that individuals are more likely to use an online dating service if they use the Internet for a greater number of tasks, and less likely to use such a service if they are trusting of others.
Attitudes towards online dating improved visibly between 2005 and 2015, the Pew Research Center found. In particular, the number of people who thought that online dating was a good way to meet people rose from 44% in 2005 to 59%. Although only a negligible number of people dated online in 2005, that rose to 11% in 2013 and then 15% in 2015. In particular, the number of American adults who had used an online dating site went from 9% in 2013 to 12% in 2015 while those who used an online dating software application on their mobile phones jumped from 3% to 9% during the same period. This increase was driven mainly by people aged 18 to 24, for whom usage almost tripled. At the same time, usage among those between the ages of 55 and 64 doubled.
According to a 2015 study by the Pew Research Center, people who had used online dating services had a higher opinion of such services than those who had not. 80% of the users said that online dating sites are a good way to meet potential partners.
In 2016, Consumer Reports surveyed approximately 115,000 online dating service subscribers across multiple platforms and found that while 44 percent of survey respondents stated that usage of online dating services led to a serious long-term intimate relationship or marriage, a subset of approximately 9,600 subscribers that had used at least one online dating service within the previous two years rated satisfaction with the services they used lower than Consumer Reports surveys of consumer satisfaction with technical support services and rated satisfaction with free online dating services as slightly more satisfactory than services with paid subscriptions. In the October 2019 Pew Research Center survey, 57% of survey respondents who had used online dating said their experiences on the platforms was very or somewhat positive while 42% said their experiences were very or somewhat negative, and 76% of survey respondents felt that online dating has had neither a positive or negative effect on dating and relationships or a mostly negative effect while 22% felt that online dating has had a mostly positive effect.
In a July 2022 survey of 6,034 U.S. adults conducted by the Pew Research Center, 53% of survey respondents who had used online dating said their experiences on the platforms were either very or somewhat positive while 46% said their experiences were either very or somewhat negative, 54% of all survey respondents said they believed that dating apps either made no difference in finding a partner or spouse or made doing so harder while 42% said they believed that dating apps made finding a partner or spouse easier, and 80% of survey respondents felt that online dating has had neither a positive or negative effect on dating and relationships or a mostly negative effect while 18% felt that online dating has had a mostly positive effect.
Trust and safety issues
As online dating services are not required to routinely conduct background checks on members, it is possible for profile information to be misrepresented or falsified. Also, there may be users on dating services that have illicit intentions (i.e. date rape, procurement, etc).
OKCupid once introduced a real name policy, but that was later taken removed due to unpopularity with its users.
Only some online dating services are providing important safety information such as STD status of its users or other infectious diseases, but many do not.
Some online dating services which are popular amongst members of queer communities are sometimes used by people as a means of meeting these audiences for the purpose of gaybashing or trans bashing.
A form of misrepresentation is that members may lie about their height, weight, age, or marital status in an attempt to market or brand themselves in a particular way. Users may also carefully manipulate profiles as a form of impression management. Online daters have raised concerns about ghosting, the practice of ceasing all communication with a person without explaining why. Ghosting appears to be becoming more common. Various explanations have been suggested, but social media is often blamed, as are dating apps and the relative anonymity and isolation in modern-day dating and hookup culture, which make it easier to behave poorly with few social repercussions.
Online dating site members may try to balance an accurate representation with maintaining their image in a desirable way. One study found that nine out of ten participants had lied on at least one attribute, though lies were often slight; weight was the most lied about attribute, and age was the least lied about. Furthermore, knowing a large amount of superficial information about a potential partner's interests may lead to a false sense of security when meeting up with a new person. Gross misrepresentation may be less likely on matrimonials sites than on casual dating sites.
Some profiles may not even represent real humans but rather they may be fake "bait profiles" placed online by site owners to attract new paying members, or "spam profiles" created by advertisers to market services and products.
Opinions on regarding the safety of online dating are mixed. Over 50% of research participants in a 2011 study did not view online dating as a dangerous activity, whereas 43% thought that online dating involved risk. Date rape is a form of acquaintance rape and dating violence. The two phrases are often used interchangeably, but date rape specifically refers to a rape in which there has been some sort of romantic or potentially sexual relationship between the two parties. Acquaintance rape also includes rapes in which the victim and perpetrator have been in a non-romantic, non-sexual relationship, for example as co-workers or neighbors. According to the United States Bureau of Justice Statistics (BJS), date rapes are among the most common forms of rape cases. Date rape most commonly takes place among college students when alcohol is involved or date rape drugs are taken. One of the most targeted groups are women between the ages of 16 and 24.
In the October 2019 Pew Research Center survey, 53% of survey respondents said believed that dating apps were a very or somewhat safe way to meet potential partners while 46% believed they were a not too safe or not at all safe way to do so, and 50% online dating respondents said that they believed that scam accounts were common. In the July 2022 Pew Research Center survey, 49% of survey respondents said believed that dating apps were a not too safe or not at all safe way to meet potential partners while 48% believed they were a very or somewhat safe way to do so, and 52% online dating respondents said that they believed that scam accounts were common.
In response to these issues, over 120 Facebook groups named Are We Dating The Same Guy? were created where women share red flags about men and check that he is not dating another person. It is done by taking screenshots of a man's dating profile and posting it onto her city's designated Facebook group, asking "any tea?". Other users in the group will then share information about the man and share warnings. The groups are moderated by volunteers, and have been described as a feminist group.
Billing complaints
Online subscription-based services can suffer from complaints about billing practices. Some online dating service providers may have fraudulent membership fees or credit card charges. Some sites do not allow members to preview available profiles before paying a subscription fee. Furthermore, different functionalities may be offered to members who have paid or not paid for subscriptions, resulting in some confusion around who can view or contact whom.
Consolidation within the online dating industry has led to different newspapers and magazines now advertising the same website database under different names. In the UK, for example, Time Out ("London Dating"), The Times ("Encounters"), and The Daily Telegraph ("Kindred Spirits"), all offer differently named portals to the same service—meaning that a person who subscribes through more than one publication has unwittingly paid more than once for access to the same service.
Imbalanced gender ratios
Little is known about the sex ratio controlled for age. eHarmony's membership is about 57% female and 43% male, whereas the ratio at Match.com is about the reverse of that. On specialty niche websites where the primary demographic is male, there is typically a very unbalanced ratio of male to female or female to male. As of June 2015, 62% of Tinder users were male and 38% were female.
Studies have suggested that men are far more likely to send messages on dating sites than women. In addition, men tend to message the most attractive women regardless of their own attractiveness. This leads to the most attractive women on these sites receiving an overwhelming number of messages, which can in some cases result in them leaving the site.
There is some evidence that there may be differences in how women online rate male attractiveness as opposed to how men rate female attractiveness. The distribution of ratings given by men of female attractiveness appears to be the normal distribution, while ratings of men given by women is highly skewed, with 80% of men rated as below average.
Allegations of discrimination
Gay rights groups have complained that certain websites that restrict their dating services to heterosexual couples are discriminating against homosexuals. Homosexual customers of the popular eHarmony dating website have made many attempts to litigate discriminatory practices. eHarmony was sued in 2007 by a lesbian claiming that "[s]uch outright discrimination is hurtful and disappointing for a business open to the public in this day and age." In light of discrimination by sexual orientation by dating websites, some services such as GayDar.net and Chemistry.com cater more to homosexual dating.
Lawsuits filed against online dating services
A 2011 class action lawsuit alleged Match.com failed to remove inactive profiles, did not accurately disclose the number of active members, and does not police its site for fake profiles; the inclusion of expired and spam profiles as valid served to both artificially inflate the total number of profiles and camouflage a skewed gender ratio in which active users were disproportionately single males. The suit claimed up to 60 percent were inactive profiles, fake or fraudulent users. Some of the spam profiles were alleged to be using images of porn actresses, models, or people from other dating sites. Former employees alleged Match routinely and intentionally over-represented the number of active members on the website and a huge percentage were not real members but 'filler profiles'.
A 2012 class action against Successful Match ended with a November 2014 California jury award of $1.4 million in compensatory damages and $15 million in punitive damages. SuccessfulMatch operated a dating site for people with STDs, PositiveSingles, which it advertised as offering a "fully anonymous profile" which is "100% confidential". The company failed to disclose that it was placing those same profiles on a long list of affiliate site domains such as GayPozDating.com, AIDSDate.com, HerpesInMouth.com, ChristianSafeHaven.com, MeetBlackPOZ.com, HIVGayMen.com, STDHookup.com, BlackPoz.com, and PositivelyKinky.com. This falsely implied that those users were black, Christian, gay, HIV-positive or members of other groups with which the registered members did not identify. The jury found PositiveSingles guilty of fraud, malice, and oppression as the plaintiffs' race, sexual orientation, HIV status, and religion were misrepresented by exporting each dating profile to niche sites associated with each trait.
In 2013, a former employee sued adultery website Ashley Madison claiming repetitive strain injuries as creating 1000 fake profiles in one three week span "required an enormous amount of keyboarding" which caused the worker to develop severe pain in her wrists and forearms. AshleyMadison's parent company, Avid Life Media, countersued in 2014, alleging the worker kept confidential documents, including copies of her "work product and training materials". The firm claimed the fake profiles were for "quality assurance testing" to test a new Brazilian version of the site for "consistency and reliability".
In January 2014, an already-married Facebook user attempting to close a pop-up advertisement for Zoosk.com found that one click instead copied personal info from her Facebook profile to create an unwanted online profile seeking a mate, leading to a flood of unexpected responses from amorous single males.
In 2014, It's Just Lunch International was the target of a New York class action alleging unjust enrichment as IJL staff relied on a uniform, misleading script which informed prospective customers during initial interviews that IJL already had at least two matches in mind for those customers' first dates regardless of whether or not that was true.
In 2014, the US Federal Trade Commission fined UK-based JDI Dating (a group of 18 websites, including Cupidswand.com and FlirtCrowd.com) over US$600000, finding that "the defendants offered a free plan that allowed users to set up a profile with personal information and photos. As soon as a new user set up a free profile, he or she began to receive messages that appeared to be from other members living nearby, expressing romantic interest or a desire to meet. However, users were unable to respond to these messages without upgrading to a paid membership ... [t]he messages were almost always from fake, computer-generated profiles — 'Virtual Cupids' — created by the defendants, with photos and information designed to closely mimic the profiles of real people." The FTC also found that paid memberships were being renewed without client authorisation.
On June 30, 2014, co-founder and former marketing vice president of Tinder, Whitney Wolfe, filed a sexual harassment and sex discrimination suit in Los Angeles County Superior Court against IAC-owned Match Group, the parent company of Tinder. The lawsuit alleged that her fellow executives and co-founders Rad and Mateen had engaged in discrimination, sexual harassment, and retaliation against her, while Tinder's corporate supervisor, IAC's Sam Yagan, did nothing. IAC suspended CMO Mateen from his position pending an ongoing investigation, and stated that it "acknowledges that Mateen sent private messages containing 'inappropriate content,' but it believes Mateen, Rad and the company are innocent of the allegations". In December 2018, The Verge reported that Tinder had dismissed Rosette Pambakian, the company's vice president of marketing and communication who had accused Tinder's former CEO Greg Blatt of sexual assault, along with several other employees who were part of the group of Tinder employees who had previously sued the Match Group for $2 billion.
Government regulation
U.S. government regulation of dating services began with the International Marriage Broker Regulation Act (IMBRA) which took effect in March 2007 after a federal judge in Georgia upheld a challenge from the dating site European Connections. The law requires dating services meeting specific criteria—including having as their primary business to connect U.S. citizens/residents with foreign nationals—to conduct, among other procedures, sex offender checks on U.S. customers before contact details can be provided to the non-U.S. citizen. In 2008, the state of New Jersey passed a law which requires the sites to disclose whether they perform background checks.
In the People's Republic of China, using a transnational matchmaking agency involving a monetary transaction is illegal. The Philippines prohibits the business of organizing or facilitating marriages between Filipinas and foreign men under the Republic Act 6955 (the Anti-Mail-Order Bride Law) of June 13, 1990; this law is routinely circumvented by basing mail-order bride websites outside the country.
Singapore's Social Development Network is the governmental organization facilitating dating activities in the country. Singapore's government has actively acted as a matchmaker for singles for the past few decades, and thus only 4% of Singaporeans have ever used an online dating service, despite the country's high rate of internet penetration.
In December 2010, a New York State Law called the "Internet Dating Safety Act" (S5180-A) went into effect that requires online dating sites with customers in New York State to warn users not to disclose personal information to people they do not know.
| Technology | Internet | null |
169720 | https://en.wikipedia.org/wiki/Melon | Melon | A melon is any of various plants of the family Cucurbitaceae with sweet, edible, and fleshy fruit. It can also specifically refer to Cucumis melo, commonly known as the "true melon" or simply "melon". The term "melon" can apply to both the plant and its fruit. Botanically, a melon is a kind of berry, specifically a "pepo". The word melon derives from Latin , which is the latinization of the Greek (mēlopepōn), meaning "melon", itself a compound of (mēlon), "apple", treefruit (of any kind)" and (pepōn), amongst others "a kind of gourd or melon". Many different cultivars have been produced, particularly of the true melon, such as the cantaloupe and honeydew.
History
Melons were thought to have originated in Africa. However, recent studies suggest a Southwest Asian origin, especially Iran and India; from there, they gradually began to appear in Europe toward the end of the Western Roman Empire. Melons are known to have been grown by the ancient Egyptians. However, recent discoveries of melon seeds dated between 1350 and 1120 BCE in Nuragic sacred wells have shown that melons were first brought to Europe by the Nuragic civilization of Sardinia during the Bronze Age. Melons were among the earliest plants to be domesticated in the Old World and among the first crop species brought by westerners to the New World. Early European settlers in the New World are recorded as growing honeydew and casaba melons as early as the 1600s. A number of Native American tribes in New Mexico, including Acoma, Cochiti, Isleta, Navajo, Santo Domingo and San Felipe, maintain a tradition of growing their own characteristic melon cultivars, derived from melons originally introduced by the Spanish. Organizations like Native Seeds/SEARCH have made an effort to collect and preserve these and other heritage seeds.
Melons by genus
Benincasa
Winter melon (B. hispida) is the only member of the genus Benincasa. The mature winter melon is a cooking vegetable that is widely used in Asia, especially in India. The immature melons are used as a culinary fruit (e.g., to make a distinctive fruit drink).
Citrullus
Egusi (C. lanatus) is a wild melon, similar in appearance to the watermelon. The flesh is inedible, but the seeds are a valuable food source in Africa. Other species that have the same culinary role, and that are also called egusi include Melothria sphaerocarpa (syn. Cucumeropsis mannii) and Lagenaria siceraria.
Watermelon (C. lanatus) originated in Africa, where evidence indicates that it has been cultivated for over 4,000 years. It is a popular summer fruit in all parts of the world.
Cucumis
Melons in genus Cucumis are culinary fruits, and include the majority of culinary melons. All but a handful of culinary melon varieties belong to the species Cucumis melo L.
Horned melon (C. metuliferus), a traditional food plant in Africa with distinctive spikes. Now grown in California, Chile, Australia and New Zealand as well.
True melon (C. melo)
C. melo cantalupensis, with skin that is rough and warty, not netted.
The European cantaloupe, with lightly ribbed, pale green skin, was domesticated in the 18th century, in Cantalupo in Sabina, Italy, by the pope's gardener. It is also known as a 'rockmelon' in Australia and New Zealand. Varieties include the French Charentais and the Burpee Seeds hybrid Netted Gem, introduced in the 19th century. The Yubari King is a highly prized Japanese cantaloupe cultivar.
The Persian melon resemble a large cantaloupe with a darker green rind and a finer netting.
C. melo inodorus, casabas, honeydew, and Asian melons
Argos, a large, oblong, with orange wrinkled skin, orange flesh, strong aroma. A characteristic is its pointed ends. Growing in some areas of Greece, from which it gets its name.
Banana melon, an heirloom variety with salmon-colored flesh and an elongated banana shape and yellow rind
Canary melon, a large, bright-yellow melon with a pale green to white inner flesh.
Casaba, bright yellow, with a smooth, furrowed skin. Less flavorful than other melons, but keeps longer.
Crenshaw melon, a hybrid between a Casaba melon and a Persian melon that is described to have a very sweet flavor
Gaya melon, originally from Japan, a honeydew cultivar that is ivory in color and has a mild, sweet flavor
Hami melon, originally from Hami, Xinjiang, China. Flesh is sweet and crisp.
Honeydew, with a sweet, juicy, green-colored flesh. Grown as bailan melon in Lanzhou, China. There is a second variety which has yellow skin, white flesh and tastes like a moist pear.
Honeymoon melon, a variety of honeydew with golden rind and bright green flesh and a sweet flavor
Kajari melon, a sweet honeydew cultivar that is red-orange in color with green stripes reminiscent of a beach ball
Kolkhoznitsa melon, with smooth, yellow skin and dense, white flesh.
Japanese melons (including the Sprite melon).
Korean melon, a yellow melon with white lines running across the fruit and white inside. Can be crisp and slightly sweet or juicy when left to ripen longer.
Mirza melon, a large, cream-colored melon native to Central Asia with a sweet, savory flavor
Oriental pickling melon
Pixie melon, a sweet, palm-sized cantaloupe cultivar with a strange, cracked-looking netting
() or Santa Claus melon, a melon with a blotchy green skin and white sweet-tasting flesh.
Sugar melon, a smooth, white, round fruit.
Tiger melon, an orange, yellow and black striped melon from Turkey with a soft pulp.
C. melo reticulatus, true muskmelons, with netted (reticulated) skin.
North American cantaloupe, distinct from the European cantaloupe, with the net-like skin pattern common to other C. melo reticulatus varieties.
Galia (or Ogen), small and very juicy with either faint green or rosy pink flesh.
Sharlyn melons, with taste between honeydew and cantaloupes, netted skin, greenish-orange rind, and white flesh.
C. melo agrestis, Wilder melon cultivars, with smooth skin, and tart or bland taste. Often confused with cucumbers (Dosakai, Lemon Cucumber, Pie Melons).
C. melo conomon, Conomon Melons, Pickling Melons, with smooth skin, and ranging from tart or bland taste (pickling melon) to mild sweetness in Korean Melon.Oriental Pickling melon, Korean Melon. Closely related to wilder melons (C Melo Var Agrestis).
Modern crossbred varieties, e.g. Crenshaw (Casaba × Persian), Crane (Japanese × N.A. cantaloupe).
Gallery
Production
In 2018, world production of melons was 27 million tonnes, led by China with 46% of the total (table). Turkey, Iran, and India each produced more than 1 million tonnes.
| Biology and health sciences | Melons | Plants |
169734 | https://en.wikipedia.org/wiki/Spinach | Spinach | Spinach (Spinacia oleracea) is a leafy green flowering plant native to central and Western Asia. It is of the order Caryophyllales, family Amaranthaceae, subfamily Chenopodioideae. Its leaves are a common edible vegetable consumed either fresh, or after storage using preservation techniques by canning, freezing, or dehydration. It may be eaten cooked or raw, and the taste differs considerably; the high oxalate content may be reduced by steaming.
It is an annual plant (rarely biennial), growing as tall as . Spinach may overwinter in temperate regions. The leaves are alternate, simple, ovate to triangular, and very variable in size: long and broad, with larger leaves at the base of the plant and small leaves higher on the flowering stem. The flowers are inconspicuous, yellow-green, in diameter, and mature into a small, hard, dry, lumpy fruit cluster across containing several seeds.
In 2022, world production of spinach was 33 million tonnes, with China alone accounting for 93% of the total.
Etymology
Originally from Persian aspānāḵ, the name entered European languages from medieval Latin spinagium, which borrowed it from Arabic isbanakh. The English word "spinach" dates to the late 14th century from the Old French word espinache.
Taxonomy
Common spinach (S. oleracea) was long considered to be in the family Chenopodiaceae, but in 2003 that family was merged into the Amaranthaceae in the order Caryophyllales. Within the family Amaranthaceae sensu lato, Spinach belongs to the subfamily Chenopodioideae.
Description
As opposed to most flowering plants used as vegetables, spinach is a dioecious plant, meaning different plants can have either female or male flowers.
The flowers are small, green and wind pollinated.
History
Spinach is thought to have originated about 2,000 years ago in ancient Persia from which it was introduced to India and later to ancient China via Nepal in 647 CE as the "Persian vegetable". In 827 CE, the Arabs introduced spinach to Sicily. The first written evidence of spinach in the Mediterranean was recorded in three 10th-century works: a medical work by al-Rāzī (known as Rhazes in the West) and in two agricultural treatises, one by Ibn Waḥshīyah and the other by Qusṭus al-Rūmī. Spinach became a popular vegetable in the Arab Mediterranean and arrived in the Iberian Peninsula by the latter part of the 12th century, where Ibn al-ʻAwwām called it , 'the chieftain of leafy greens'. Spinach was also the subject of a special treatise in the 11th century by Ibn Ḥajjāj.
Spinach first appeared in England and France in the 14th century, probably via Iberia, and gained common use because it appeared in early spring when fresh local vegetables were not available. Spinach is mentioned in the first known English cookbook, the Forme of Cury (1390), where it is referred to as 'spinnedge' and 'spynoches'. During World War I, wine fortified with spinach juice was given to injured French soldiers with the intent to curtail their bleeding.
Culinary use
Nutrients
Raw spinach is 91% water, 4% carbohydrates, 3% protein, and contains negligible fat (table). In a reference serving providing of food energy, spinach has a high nutritional value, especially when fresh, frozen, steamed, or quickly boiled. It is a rich source (20% or more of the Daily Value, DV) of vitamin A, vitamin C, manganese, and folate (31-52% DV), with an especially high content of vitamin K (403% DV) (table). Spinach is a moderate source (10–19% of DV) of the B vitamins, riboflavin and vitamin B6, vitamin E, potassium, iron, magnesium, and dietary fiber (table).
Although spinach contains moderate amounts of iron and calcium, it also contains oxalates, which may inhibit absorption of calcium and iron in the stomach and small intestine. Cooked spinach has lower levels of oxalates, and its nutrients may be absorbed more completely.
Cooking spinach significantly decreases its vitamin C concentration, as vitamin C is degraded by heating. Folate levels may also be decreased, as folate tends to leach into cooking liquid.
Spinach is rich in nitrates and nitrites, which may exceed safe levels if spinach is over-consumed.
Cuisine
Spinach is eaten raw, in salads, and cooked in soups, curries, or casseroles. Dishes with spinach as a main ingredient include spinach salad, spinach soup, spinach dip, saag paneer, pkhali, ispanakhi matsvnit, and spanakopita.
In classical French cuisine, a spinach-based dish may be described as à la Florentine.
Production
In 2022, world production of spinach was 33 million tonnes, with China alone accounting for 93% of the total.
Marketing and safety
Fresh spinach is sold loose, bunched, or packaged fresh in bags. Fresh spinach loses much of its nutritional value with storage of more than a few days. Fresh spinach is packaged in air, or in nitrogen gas to extend shelf life. While refrigeration slows this effect to about eight days, fresh spinach loses most of its folate and carotenoid content over this period of time. For longer storage, it is canned, or blanched or cooked and frozen.
Some packaged spinach is exposed to radiation to kill any harmful bacteria. The Food and Drug Administration approves of irradiation of spinach leaves up to 4.0 kilograys, having no or only a minor effect on nutrient content.
Spinach may be high in cadmium contamination depending on the soil and location where the spinach is grown.
Due to spinach's high content of vitamin K, individuals taking the anticoagulant warfarin, which acts by inhibiting vitamin K, are instructed to minimize consumption of spinach (and other dark green leafy vegetables).
In popular culture
The comics and cartoon character Popeye the Sailor Man is portrayed as gaining strength by consuming canned spinach. The accompanying song lyric is: "I'm strong to the finich , 'cuz I eats me spinach." This is usually attributed to the iron content of spinach, but in a 1932 strip, Popeye states that "spinach is full of vitamin A" and that is what makes people strong and healthy. As it happens, spinach is not a better source of dietary iron than many other vegetables. The false idea that spinach is an especially good source of dietary iron is an academic urban legend.
| Biology and health sciences | Caryophyllales | null |
169785 | https://en.wikipedia.org/wiki/Dewberry | Dewberry | The dewberries are a group of species in the genus Rubus, section Rubus, closely related to the blackberries. They are small trailing (rather than upright or high-arching) brambles with aggregate fruits, reminiscent of the raspberry, but are usually purple to black instead of red.
Description
The plants do not have upright canes like some other Rubus species, but have stems that trail along the ground, putting forth new roots along the length of the stem. The stems are covered with fine spines or stickers. Around March and April, the plants start to grow white flowers that develop into small green berries. The tiny green berries grow red and then a deep purple-blue as they ripen. When the berries are ripe, they are tender and difficult to pick in any quantity without squashing them. The berries are sweet and often less seedy than blackberries.
In the winter the leaves often remain on the stems, but may turn dark red.
The European dewberry, Rubus caesius, grows more upright like other brambles. Its fruits are a deep, almost black, purple and are coated with a thin layer or 'dew' of waxy droplets. Thus, they appear sky-blue (caesius being Latin for pale blue). Its fruits are small and retain a markedly tart taste even when fully ripe.
Species
Rubus Section Caesii, European dewberry
European dewberry, Rubus caesius L.
Rubus Section Flagellares, American dewberries
Rubus aboriginum Rydb., synonyms:
Rubus almus (L.H. Bailey) L.H.Bailey
Rubus austrinus L.H.Bailey
Rubus bollianus L.H.Bailey
Rubus clair-brownii''' L.H.Bailey
Rubus decor L.H. Bailey
Rubus flagellaris Willd. var. almus L.H.Bailey
Rubus foliaceus L.H. Bailey
Rubus ignarus L.H. Bailey
Rubus ricei L.H. Bailey
Aberdeen dewberry, Rubus depavitus L.H.Bailey
Northern dewberry, Rubus flagellaris Willd.
Swamp dewberry, Rubus hispidus L.
Upland dewberry, Rubus invisus (L.H.Bailey) Britton
Pacific dewberry, Rubus ursinus Cham. & Schltdl.
Southern dewberry Rubus trivialis L.H.Bailey
Distribution and habitat
Dewberries are common throughout most of the Northern Hemisphere and are thought of as a beneficial weed. Rubus caesius is frequently restricted to coastal communities, especially sand dune systems.
Ecology
The leaves are sometimes eaten by the larvae of some Lepidoptera species including peach blossom moths.
Uses
The leaves can be used to make a herbal tea, and the berries are edible and taste sweet. They can be eaten raw, or used to make cobbler, jam, or pie.
In the late 19th and early 20th centuries, the town of Cameron, North Carolina, was known as the "dewberry capital of the world" for large scale cultivation of this berry which was shipped out for widespread consumption. Local growers made extensive use of the railroads in the area to ship them nationally and internationally.
| Biology and health sciences | Berries | Plants |
169799 | https://en.wikipedia.org/wiki/Neon%20tetra | Neon tetra | The neon tetra (Paracheirodon innesi) is a freshwater fish of the characin family (family Characidae) of order Characiformes. The type species of its genus, it is native to blackwater and clearwater streams in the Amazon basin of South America. Its bright colouring makes the fish visible to conspecifics in the dark blackwater streams, and is also the main reason for its popularity among freshwater fish hobbyists, with neon tetras being one of the most widely kept tropical fish in the world.
Range and habitat
The neon tetra is found in the western and northern Amazon basin in southeastern Colombia, eastern Peru, and western Brazil. It lives in waters with a temperature between and pH 4–7.5. It has a preference for acidic blackwater streams, but also occurs in transparent clearwater streams. It is not found in the whitewater rivers. UN FAO considers P. innesi an introduced species in Singapore and the United States. FAO considers its introduction to Singapore to be ecologically and socioeconomically beneficial, but it is not established there.
Description
The neon tetra has a light-blue back over a silver-white abdomen. The fish is characterized by an iridescent blue horizontal stripe along each side of the fish from its nose to the base of the adipose fin, and an iridescent red stripe that begins at the middle of the body and extends posteriorly to the base of the caudal fin. The fish is partially transparent (including fins) except for these markings. Sexual dimorphism is slight, with the female having a slightly larger belly, and a bent iridescent stripe rather than the male's straight stripe. During the night, the blue and red become gray or black as the fish rests. It reactivates once it becomes active in the morning. This peculiar change is due to the neon tetras capacity to change the color of its iridescent stripe in response to lighting conditions. In a light-adapted state it is blue-green, likewise in a dark-adapted state its color changes to indigo. This change is produced by guanine crystals in their cells that reflect light. The neon tetra grows to approximately in overall length.
Economics
The neon tetra was first imported from South America and was described by renowned ichthyologist George S. Myers in 1936, and named after William T. Innes. P. innesi is one of the most popular aquarium fish, with about 2 million sold in the US each month. Most neon tetras available in the United States are imported from Southeast Asia, where they are farm-raised, or to a lesser extent from Colombia, Peru, and Brazil, where they are collected from the wild. During a single month, an average of 1.8 million neon tetras with an estimated value of $175,000 are imported into the United States for the aquarium trade. With the exception of home aquarists and a few commercial farms that breed neon tetras experimentally, captive breeding on a commercial scale is nonexistent in the USA.
In the aquarium
In the wild they inhabit very soft, acidic waters (pH 4.0 to 4.8) Ideal pH for aquarium is 7.0, but a range of 6.0 to 8.0 is tolerable. They can have a lifespan of as long as ten years, but normally just two to three years in an aquarium.
Neon tetras are considered easy to keep in an aquarium of at least with a temperature range of between , a water pH of between 6.0 and 7.0, GH of below 10 dGH and KH of 1–2 dKH, and under 20 ppm of nitrate.
They are shoaling fish and must be kept in groups of at least six, but will be more active in groups of eight to 12 or more. Neon tetras are best kept in a densely planted tank to resemble their native Amazon environments.
Nutrition
Neon tetras are omnivores and will accept most flake foods, if sufficiently small, but should also have some small foods such as brine shrimp, daphnia, freeze-dried bloodworms, tubifex, which can be stuck to the side of the aquarium, and micropellet food to supplement their diets. A tropical sinking pellet is ideal, as most brands of these include natural color enhancers that bring out the color in neon tetras. Some frozen foods, including frozen blood worms, add variety to their diets.
Breeding
The male is slender, and the blue line is straighter. The female is rounder, producing a bent blue line. Some aquarists say the females look plumper when viewed from above. However, the straightness of the line and the plumpness of the female might occasionally be due to the eggs she is carrying. A neon tetra can appear slightly plump in the belly due to having overeaten.
Neon tetras need dim lighting, a DH less than one, about 5.5 pH, and a temperature of to breed. There also needs to be a lot of tannins in the water. Neon tetras are old enough to breed at 12 weeks.
Breeding neon tetras is considered to be difficult in home aquariums. However, it is becoming more common, with less than 5% of specimens currently sold in America caught in the wild, and more than 1.5 million specimens imported to America each month from fish farms.
Disease
Neon tetras are occasionally afflicted by the so-called "neon tetra disease" (NTD) or pleistophora disease, a sporozoan disease caused by Pleistophora hyphessobryconis. Despite being a well-known condition, it is generally incurable and often fatal to the fish. However this disease is also generally preventable. The disease cycle begins when microsporidian parasite spores enter the fish after it consumes infected material, such as the bodies of a dead fish, or live food such as tubifex, which may serve as intermediate hosts. The disease is most likely passed by newly acquired fish that have not been quarantined.
Symptoms include restlessness, loss of coloration, lumps on the body as cysts develop, difficulty swimming, curved spines as the disease progresses, and secondary infections, such as fin rot and bloating.
A so-called "false neon disease", which is bacterial, shows very similar symptoms. It is impossible for the home aquarist to determine for certain the difference between NTD and false NTD on the basis of visible symptoms alone, without laboratory backup. This disease has also been confused with columnaris (mouth rot, mouth fungus, 'flex').
Generally the best 'treatment' is the immediate removal of diseased fish to preserve the remaining fish, although some occasional successful treatments have been performed that include fish baths and a "medication cocktail". The use of a diatom filter, which can reduce the number of free parasites in the water, may help. As with Pleistophora neon tetra disease, prevention is most important and this disease is rare when good preventive measures are performed.
Related species
The green neon tetra (P. simulans) and black neon tetra (Hyphessobrycon herbertaxelrodi) are distinct species—the latter belongs to a different genus—and not color varieties. The cardinal tetra (P. axelrodi) is also a similar species, but its greater extent of red coloring distinguishes it from the neon tetra. The name Hyphessobrycon innesi is an obsolete synonym for P. innesi.
| Biology and health sciences | Characiformes | Animals |
169889 | https://en.wikipedia.org/wiki/Paleoproterozoic | Paleoproterozoic | The Paleoproterozoic Era (also spelled Palaeoproterozoic) is the first of the three sub-divisions (eras) of the Proterozoic eon, and also the longest era of the Earth's geological history, spanning from (2.5–1.6 Ga). It is further subdivided into four geologic periods, namely the Siderian, Rhyacian, Orosirian and Statherian.
Paleontological evidence suggests that the Earth's rotational rate ~1.8 billion years ago equated to 20-hour days, implying a total of ~450 days per year. It was during this era that the continents first stabilized.
Atmosphere
The Earth's atmosphere was originally a weakly reducing atmosphere consisting largely of nitrogen, methane, ammonia, carbon dioxide and inert gases, in total comparable to Titan's atmosphere. When oxygenic photosynthesis evolved in cyanobacteria during the Mesoarchean, the increasing amount of byproduct dioxygen began to deplete the reductants in the ocean, land surface and the atmosphere. Eventually all surface reductants (particularly ferrous iron, sulfur and atmospheric methane) were exhausted, and the atmospheric free oxygen levels soared permanently during the Siderian and Rhyacian periods in an aerochemical event called the Great Oxidation Event, which brought atmospheric oxygen from near none to up to 10% of the modern level.
Life
At the beginning of the preceding Archean eon, almost all existing lifeforms were single-cell prokaryotic anaerobic organisms whose metabolism was based on a form of cellular respiration that did not require oxygen, and autotrophs were either chemosynthetic or relied upon anoxygenic photosynthesis. After the Great Oxygenation Event, the then mainly archaea-dominated anaerobic microbial mats were devastated as free oxygen is highly reactive and biologically toxic to cellular structures. This was compounded by a 300-million-year-long global icehouse event known as the Huronian glaciation — at least partly due to the depletion of atmospheric methane, a powerful greenhouse gas — resulted in what is widely considered one of the first and most significant mass extinctions on Earth. The organisms that thrived after the extinction were mainly aerobes that evolved bioactive antioxidants and eventually aerobic respiration, and surviving anaerobes were forced to live symbiotically alongside aerobes in hybrid colonies, which enabled the evolution of mitochondria in eukaryotic organisms.
The Palaeoproterozoic represents the era from which the oldest cyanobacterial fossils, those of Eoentophysalis belcherensis from the Kasegalik Formation in the Belcher Islands of Nunavut, are known. By 1.75 Ga, thylakoid-bearing cyanobacteria had evolved, as evidenced by fossils from the McDermott Formation of Australia.
Many crown node eukaryotes (from which the modern-day eukaryotic lineages would have arisen) have been approximately dated to around the time of the Paleoproterozoic Era.
While there is some debate as to the exact time at which eukaryotes evolved,
current understanding places it somewhere in this era. Statherian fossils from the Changcheng Group in North China provide evidence that eukaryotic life was already diverse by the late Palaeoproterozoic.
Geological events
During this era, the earliest global-scale continent-continent collision belts developed. The associated continent and mountain building events are represented by the 2.1–2.0 Ga Trans-Amazonian and Eburnean orogens in South America and West Africa; the ~2.0 Ga Limpopo Belt in southern Africa; the 1.9–1.8 Ga Trans-Hudson, Penokean, Taltson–Thelon, Wopmay, Ungava and Torngat orogens in North America, the 1.9–1.8 Ga Nagssugtoqidian Orogen in Greenland; the 1.9–1.8 Ga Kola–Karelia, Svecofennian, Volhyn-Central Russian, and Pachelma orogens in Baltica (Eastern Europe); the 1.9–1.8 Ga Akitkan Orogen in Siberia; the ~1.95 Ga Khondalite Belt; the ~1.85 Ga Trans-North China Orogen in North China; and the 1.8-1.6 Ga Yavapai and Mazatzal orogenies in southern North America.
That pattern of collision belts supports the formation of a Proterozoic supercontinent named Columbia or Nuna. That continental collisions suddenly led to mountain building at large scale is interpreted as having resulted from increased biomass and carbon burial during and after the Great Oxidation Event: Subducted carbonaceous sediments are hypothesized to have lubricated compressive deformation and led to crustal thickening.
Felsic volcanism in what is now northern Sweden led to the formation of the Kiruna and Arvidsjaur porphyries.
The lithospheric mantle of Patagonia's oldest blocks formed.
| Physical sciences | Geological timescale | Earth science |
169942 | https://en.wikipedia.org/wiki/Troy%20weight | Troy weight | Troy weight is a system of units of mass that originated in the Kingdom of England in the 15th century and is primarily used in the precious metals industry. The troy weight units are the grain, the pennyweight (24 grains), the troy ounce (20 pennyweights), and the troy pound (12 troy ounces). The troy grain is equal to the grain unit of the avoirdupois system, but the troy ounce is heavier than the avoirdupois ounce, and the troy pound is lighter than the avoirdupois pound. Legally, one troy ounce (oz t) equals exactly 31.1034768 grams.
Etymology
Troy weight is generally supposed to take its name from the French market town of Troyes where English merchants traded at least as early as the early 9th century. The name troy is first attested in 1390, describing the weight of a platter, in an account of the travels in Europe of the Earl of Derby.
Charles Moore Watson (1844–1916) proposes an alternative etymology: The Assize of Weights and Measures (also known as ), one of the statutes of uncertain date from the reign of either Henry III or Edward I, thus before 1307, specifies ""—which the Public Record Commissioners translate as "troy weight". The word refers to markets. Wright's The English Dialect Dictionary lists the word troi as meaning a balance, related to the alternate form 'tron' which also means market or the place of weighing. From this, Watson suggests that 'troy' derives from the manner of weighing by balance precious goods such as bullion or drugs; in contrast to the word 'avoirdupois' used to describe bulk goods such as corn or coal, sometimes weighed in ancient times by a kind of steelyard called the auncel.
Troy weight referred to the Tower system; the earliest reference to the modern troy weights is in 1414.
History
The origin of the troy weight system is unknown. Although the name probably comes from the Champagne fairs at Troyes, in northeastern France. English troy weights were nearly identical to the troy weight system of Bremen. (The Bremen troy ounce had a mass of 480.8 British Imperial grains.)
Many aspects of the troy weight system were indirectly derived from the Roman monetary system. Before they used coins, early Romans used bronze bars of varying weights as currency. An ("heavy bronze") weighed one pound. One twelfth of an was called an , or in English, an "ounce". Before the adoption of the metric system, many systems of troy weights were in use in various parts of Europe, among them Holland troy, Paris troy, etc. Their values varied from one another by up to several percentage points. Troy weights were first used in England in the 15th century and were made official for gold and silver in 1527. The British Imperial system of weights and measures (also known as Imperial units) was established in 1824, prior to which the troy weight system was a subset of pre-Imperial English units.
The troy ounce in modern use is essentially the same as the British Imperial troy ounce (1824–1971), adopted as an official weight standard for United States coinage by act of Congress on May 19, 1828. The British Imperial troy ounce (known more commonly simply as the imperial troy ounce) was based on, and virtually identical with, the pre-1824 British troy ounce and the pre-1707 English troy ounce. (1824 was the year the British Imperial system of weights and measures was adopted; 1707 was the year of the Act of Union which created the Kingdom of Great Britain.) Troy ounces have been used in England since the early 15th century, and the English troy ounce was officially adopted for coinage in 1527. Before that time, various sorts of troy ounces were in use on the continent.
The troy ounce and grain were also part of the apothecaries' system. This was long used in medicine, but has been largely replaced by the metric system (milligrams). The only troy weight in widespread use is the British Imperial troy ounce and its American counterpart. Both are based on a grain of 0.06479891 gram (exact, by definition), with 480 grains to a troy ounce (compared with grains for an ounce avoirdupois). The British Empire abolished the 12-ounce troy pound in the 19th century. It has been retained, though rarely used, in the American system. Larger amounts of precious metals are conventionally counted in hundreds or thousands of troy ounces, or in kilograms.
Troy ounces have been and are still often used in precious metal markets in countries that otherwise use International System of Units (SI). However, the People's Bank of China which had been using troy measurements in minting Gold Pandas since 1982 from 2016 specifies Chinese bullion coins in an integer numbers of grams.
Units of measurement
Troy pound (lb t)
The troy pound (lb t) consists of twelve troy ounces and thus is . (An avoirdupois pound is approximately 21.53% heavier at , and consists of sixteen avoirdupois ounces).
Troy ounce (oz t)
A troy ounce weighs 480 grains. Since the implementation of the international yard and pound agreement of 1 July 1959, the grain measure is defined as precisely . Thus one troy ounce = × /grain = . Since the ounce avoirdupois is defined as 437.5 grains, a troy ounce is exactly = or about 1.09714 ounces avoirdupois or about 9.7% more. The troy ounce for trading precious metals is considered to be sufficiently approximated by 31.10 g in EU directive 80/181/EEC.
The Dutch troy system is based on a mark of 8 ounces, the ounce of 20 engels (pennyweights), the engel of 32 as. The mark was rated as 3,798 troy grains or 246.084 grams. The divisions are identical to the tower system.
Pennyweight (dwt)
The pennyweight symbol is dwt. One pennyweight weighs 24 grains, and 20 pennyweights make one troy ounce. Because there were 12 troy ounces in the old troy pound, there would have been 240 pennyweights to the pound (mass) just as there were 240 pennies in the original pound-sterling. However, prior to 1526, the English pound sterling was based on the tower pound, which is of a troy pound. The d in dwt stands for denarius, the ancient Roman coin that equates loosely to a penny. The symbol d for penny can be recognized in the form of British pre-decimal pennies, in which pounds, shillings, and pence were indicated using the symbols £, s, and d, respectively.
Troy grain
There is no specific 'troy grain'. All Imperial systems use the same measure of mass called a grain (historically of barley), each weighing of an avoirdupois pound (and thus a little under 65 milligrams).
Mint masses
Mint masses, also known as moneyers' masses, were legalized by Act of Parliament dated 17 July 1649 entitled An Act touching the monies and coins of England. A grain is 20 mites, a mite is 24 droits, a droit is 20 perits, a perit is 24 blanks.
Conversions
The troy system was used in the apothecaries' system but with different further subdivisions.
| Physical sciences | Measurement systems | Basics and measurement |
169946 | https://en.wikipedia.org/wiki/Grain%20%28unit%29 | Grain (unit) | A grain is a unit of measurement of mass, and in the troy weight, avoirdupois, and apothecaries' systems, equal to exactly . It is nominally based upon the mass of a single ideal seed of a cereal. From the Bronze Age into the Renaissance, the average masses of wheat and barley grains were part of the legal definitions of units of mass. Expressions such as "thirty-two grains of wheat, taken from the middle of the ear" appear to have been ritualistic formulas. Another source states that it was defined such that 252.458 units would balance of distilled water at an ambient air-water pressure and temperature of and respectively. Another book states that Captain Henry Kater, of the British Standards Commission, arrived at this value experimentally.
The grain was the legal foundation of traditional English weight systems, and is the only unit that is equal throughout the troy, avoirdupois, and apothecaries' systems of mass. The unit was based on the weight of a single grain of barley which was equal to about the weight of a single grain of wheat. The fundamental unit of the pre-1527 English weight system, known as Tower weights, was based on the wheat grain. The tower "wheat" grain was defined as exactly (≈) of the troy "barley" grain.
Since the implementation of the international yard and pound agreement of 1 July 1959, the grain or troy grain (symbol: gr) measure has been defined in terms of units of mass in the International System of Units as precisely . One gram is thus approximately equivalent to . The unit formerly used by jewellers to measure pearls, diamonds, and other precious stones, called the jeweller's grain or pearl grain, is equal to . The grain was also the name of a traditional French unit equal to .
In both British Imperial units and United States customary units, there are precisely 7,000 grains per avoirdupois pound, and 5,760 grains per troy pound or apothecaries' pound. It is obsolete in the United Kingdom and, like most other non-SI units, it has no basis in law and cannot be used in commerce.
Current usage
Grains are commonly used to measure the mass of bullets and propellants. In archery, the grain is the standard unit used to weigh arrows.
In North America, the hardness of water is often measured in grains per U.S. gallon () of calcium carbonate equivalents. Otherwise, water hardness is measured in the dimensionless unit of parts per million (), numerically equivalent to concentration measured in milligrams per litre. One grain per U.S. gallon is approximately . Soft water contains of calcium carbonate equivalents, while hard water contains .
Though no longer recommended, in the U.S., grains are still used occasionally in medicine as part of the apothecaries' system, especially in prescriptions for older medicines such as aspirin or phenobarbital. For example, the dosage of a standard tablet of aspirin is sometimes given as . In that example the grain is approximated to , though the grain can also be approximated to , depending on the medication and manufacturer. The apothecaries' system has its own system of notation, in which the units symbol or abbreviation is followed by the quantity in lower case Roman numerals. For amounts less than one, the quantity is written as a fraction, or for one half, ss (or variations such as ss., ṡṡ, or s̅s̅). Therefore, a prescription for tablets containing 325 mg of aspirin and 30 mg of codeine can be written "ASA gr. v c̄ cod. gr. ss tablets" (using the medical abbreviations ASA for acetylsalicylic acid [aspirin], c̄ for "with", and cod. for codeine). The apothecaries' system has gradually been replaced by the metric system, and the use of the grain in prescriptions is now rare.
In the U.S., particulate emission levels, used to monitor and regulate pollution, are sometimes measured in grains per cubic foot instead of the more usual by volume. This is the same unit commonly used to measure the amount of moisture in the air, also known as the absolute humidity. The SI unit used to measure particulate emissions and absolute humidity is mg/m. One grain per cubic foot is approximately .
History
At least since antiquity, grains of wheat or barley were used by Mediterranean traders to define units of mass; along with other seeds, especially those of the carob tree. According to a longstanding tradition, one carat (the mass of a carob seed) was equivalent to the weight of four wheat grains or three barleycorns. Since the weights of these seeds are highly variable, especially that of the cereals as a function of moisture, this is a convention more than an absolute law.
The history of the modern British grain can be traced back to a royal decree in thirteenth century England, re-iterating decrees that go back as far as King Offa (eighth century). The Tower pound was one of many monetary pounds of 240 silver pennies.
The pound in question is the Tower pound. The Tower pound, abolished in 1527, consisted of 12 ounces like the troy pound, but was (≈6%) lighter. The weight of the original sterling pennies was 22½ troy grains, or 32 "Tower grains".
Physical grain weights were made and sold commercially at least as late as the early 1900s, and took various forms, from squares of sheet metal to manufactured wire shapes and coin-like weights.
The troy pound was only "the pound of Pence, Spices, Confections, as of Electuaries", as such goods might be measured by a troi or small balance. The old troy standard was set by King Offa's currency reform, and was in full use in 1284 (Assize of Weights and Measures, King Edward I), but was restricted to currency (the pound of pennies) until it was abolished in 1527. This pound was progressively replaced by a new pound, based on the weight of 120 silver dirhems of 48 grains. The new pound used a barley-corn grain, rather than a wheat grain.
Avoirdupois (goods of weight) refers to those things measured by the lesser but quicker balances: the bismar or auncel, the Roman balance, and the steelyard. The original mercantile pound of 25 shillings or 15 (Tower) ounces was displaced by, variously, the pound of the Hanseatic League (16 tower ounces) and by the pound of the then-important wool trade (16 ounces of 437 grains). A new pound of grains was inadvertently created as 16 troy ounces, referring to the new troy rather than the old troy. Eventually, the wool pound won out.
The avoirdupois pound was defined in prototype, rated as to grains. In the Imperial Weights and Measures Act 1824 (5 Geo. 4. c. 74), the avoirdupois pound was defined as grains exactly. The Weights and Measures Act 1855 authorised Miller's new standards to replace those lost in the fire that destroyed the Houses of Parliament. The standard was an avoirdupois pound, the grain being defined as of it.
The division of the carat into four grains survives in both senses well into the early twentieth century. For pearls and diamonds, weight is quoted in carats, divided into four grains. The carat was eventually set to 205 milligrams (1877), and later 200 milligrams. For touch or fineness of gold, the fraction of gold was given as a weight, the total being a solidus of 24 carats or 96 grains.
| Physical sciences | Mass and weight | Basics and measurement |
170045 | https://en.wikipedia.org/wiki/Four-stroke%20engine | Four-stroke engine | A four-stroke (also four-cycle) engine is an internal combustion (IC) engine in which the piston completes four separate strokes while turning the crankshaft. A stroke refers to the full travel of the piston along the cylinder, in either direction. The four separate strokes are termed:
Intake: Also known as induction or suction. This stroke of the piston begins at top dead center (T.D.C.) and ends at bottom dead center (B.D.C.). In this stroke the intake valve must be in the open position while the piston pulls an air-fuel mixture into the cylinder by producing a partial vacuum (negative pressure) in the cylinder through its downward motion.
Compression: This stroke begins at B.D.C, or just at the end of the suction stroke, and ends at T.D.C. In this stroke the piston compresses the air-fuel mixture in preparation for ignition during the power stroke (below). Both the intake and exhaust valves are closed during this stage.
Combustion: Also known as power or ignition. This is the start of the second revolution of the four stroke cycle. At this point the crankshaft has completed a full 360 degree revolution. While the piston is at T.D.C. (the end of the compression stroke) the compressed air-fuel mixture is ignited by a spark plug (in a gasoline engine) or by heat generated by high compression (diesel engines), forcefully returning the piston to B.D.C. This stroke produces mechanical work from the engine to turn the crankshaft.
Exhaust: Also known as outlet. During the exhaust stroke, the piston, once again, returns from B.D.C. to T.D.C. while the exhaust valve is open. This action expels the spent air-fuel mixture through the exhaust port.
Four-stroke engines are the most common internal combustion engine design for motorized land transport, being used in automobiles, trucks, diesel trains, light aircraft and motorcycles. The major alternative design is the two-stroke cycle.
History
Otto cycle
Nikolaus August Otto was a traveling salesman for a grocery concern. In his travels, he encountered the internal combustion engine built in Paris by Belgian expatriate Jean Joseph Etienne Lenoir. In 1860, Lenoir successfully created a double-acting engine that ran on illuminating gas at 4% efficiency. The 18 litre Lenoir Engine produced only 2 horsepower. The Lenoir engine ran on illuminating gas made from coal, which had been developed in Paris by Philip Lebon.
In testing a replica of the Lenoir engine in 1861, Otto became aware of the effects of compression on the fuel charge. In 1862, Otto attempted to produce an engine to improve on the poor efficiency and reliability of the Lenoir engine. He tried to create an engine that would compress the fuel mixture prior to ignition, but failed as that engine would run no more than a few minutes prior to its destruction. Many other engineers were trying to solve the problem, with no success.
In 1864, Otto and Eugen Langen founded the first internal combustion engine production company, NA Otto and Cie (NA Otto and Company). Otto and Cie succeeded in creating a successful atmospheric engine that same year.
The factory ran out of space and was moved to the town of Deutz, Germany in 1869, where the company was renamed to Deutz Gasmotorenfabrik AG (The Deutz Gas Engine Manufacturing Company). In 1872, Gottlieb Daimler was technical director and Wilhelm Maybach was the head of engine design. Daimler was a gunsmith who had worked on the Lenoir engine.
By 1876, Otto and Langen succeeded in creating the first internal combustion engine that compressed the fuel mixture prior to combustion for far higher efficiency than any engine created to this time.
Daimler and Maybach left their employ at Otto and Cie and developed the first high-speed Otto engine in 1883. In 1885, they produced the first automobile to be equipped with an Otto engine. The Daimler Reitwagen used a hot-tube ignition system and the fuel known as Ligroin to become the world's first vehicle powered by an internal combustion engine. It used a four-stroke engine based on Otto's design. The following year, Karl Benz produced a four-stroke engined automobile that is regarded as the first car.
In 1884, Otto's company, then known as Gasmotorenfabrik Deutz (GFD), developed electric ignition and the carburetor. In 1890, Daimler and Maybach formed a company known as Daimler Motoren Gesellschaft. Today, that company is Daimler-Benz.
Atkinson cycle
The Atkinson-cycle engine is a type of single stroke internal combustion engine invented by James Atkinson in 1882. The Atkinson cycle is designed to provide efficiency at the expense of power density, and is used in some modern hybrid electric applications.
The original Atkinson-cycle piston engine allowed the intake, compression, power, and exhaust strokes of the four-stroke cycle to occur in a single turn of the crankshaft and was designed to avoid infringing certain patents covering Otto-cycle engines.
Due to the unique crankshaft design of the Atkinson, its expansion ratio can differ from its compression ratio and, with a power stroke longer than its compression stroke, the engine can achieve greater thermal efficiency than a traditional piston engine. While Atkinson's original design is no more than a historical curiosity, many modern engines use unconventional valve timing to produce the effect of a shorter compression stroke/longer power stroke, thus realizing the fuel economy improvements the Atkinson cycle can provide.
Diesel cycle
The diesel engine is a technical refinement of the 1876 Otto-cycle engine. Where Otto had realized in 1861 that the efficiency of the engine could be increased by first compressing the fuel mixture prior to its ignition, Rudolf Diesel wanted to develop a more efficient type of engine that could run on much heavier fuel. The Lenoir, Otto Atmospheric, and Otto Compression engines (both 1861 and 1876) were designed to run on Illuminating Gas (coal gas). With the same motivation as Otto, Diesel wanted to create an engine that would give small industrial companies their own power source to enable them to compete against larger companies, and like Otto, to get away from the requirement to be tied to a municipal fuel supply. Like Otto, it took more than a decade to produce the high-compression engine that could self-ignite fuel sprayed into the cylinder. Diesel used an air spray combined with fuel in his first engine.
During initial development, one of the engines burst, nearly killing Diesel. He persisted, and finally created a successful engine in 1893. The high-compression engine, which ignites its fuel by the heat of compression, is now called the diesel engine, whether a four-stroke or two-stroke design.
The four-stroke diesel engine has been used in the majority of heavy-duty applications for many decades. It uses a heavy fuel containing more energy and requiring less refinement to produce. The most efficient Otto-cycle engines run near 30% thermal efficiency.
Thermodynamic analysis
The thermodynamic analysis of the actual four-stroke and two-stroke cycles is not a simple task. However, the analysis can be simplified significantly if air standard assumptions are utilized. The resulting cycle, which closely resembles the actual operating conditions, is the Otto cycle.
During normal operation of the engine, as the air/fuel mixture is being compressed, an electric spark is created to ignite the mixture. At low rpm this occurs close to TDC (Top Dead Centre). As engine rpm rises, the speed of the flame front does not change so the spark point is advanced earlier in the cycle to allow a greater proportion of the cycle for the charge to combust before the power stroke commences. This advantage is reflected in the various Otto engine designs; the atmospheric (non-compression) engine operates at 12% efficiency whereas the compressed-charge engine has an operating efficiency around 30%.
Fuel considerations
A problem with compressed charge engines is that the temperature rise of the compressed charge can cause pre-ignition. If this occurs at the wrong time and is too energetic, it can damage the engine. Different fractions of petroleum have widely varying flash points (the temperatures at which the fuel may self-ignite). This must be taken into account in engine and fuel design.
The tendency for the compressed fuel mixture to ignite early is limited by the chemical composition of the fuel. There are several grades of fuel to accommodate differing performance levels of engines. The fuel is altered to change its self ignition temperature. There are several ways to do this. As engines are designed with higher compression ratios the result is that pre-ignition is much more likely to occur since the fuel mixture is compressed to a higher temperature prior to deliberate ignition. The higher temperature more effectively evaporates fuels such as gasoline, which increases the efficiency of the compression engine. Higher compression ratios also mean that the distance that the piston can push to produce power is greater (which is called the expansion ratio).
The octane rating of a given fuel is a measure of the fuel's resistance to self-ignition. A fuel with a higher numerical octane rating allows for a higher compression ratio, which extracts more energy from the fuel and more effectively converts that energy into useful work while at the same time preventing engine damage from pre-ignition. High Octane fuel is also more expensive.
Many modern four-stroke engines employ gasoline direct injection or GDI. In a gasoline direct-injected engine, the injector nozzle protrudes into the combustion chamber. The direct fuel injector injects gasoline under a very high pressure into the cylinder during the compression stroke, when the piston is closer to the top.
Diesel engines by their nature do not have concerns with pre-ignition. They have a concern with whether or not combustion can be started. The description of how likely Diesel fuel is to ignite is called the Cetane rating. Because Diesel fuels are of low volatility, they can be very hard to start when cold. Various techniques are used to start a cold Diesel engine, the most common being the use of a glow plug.
Design and engineering principles
Power output limitations
The maximum amount of power generated by an engine is determined by the maximum amount of air ingested. The amount of power generated by a piston engine is related to its size (cylinder volume), whether it is a two-stroke engine or four-stroke design, volumetric efficiency, losses, air-to-fuel ratio, the calorific value of the fuel, oxygen content of the air and speed (RPM). The speed is ultimately limited by material strength and lubrication. Valves, pistons and connecting rods suffer severe acceleration forces. At high engine speed, physical breakage and piston ring flutter can occur, resulting in power loss or even engine destruction. Piston ring flutter occurs when the rings oscillate vertically within the piston grooves they reside in. Ring flutter compromises the seal between the ring and the cylinder wall, which causes a loss of cylinder pressure and power. If an engine spins too quickly, valve springs cannot act quickly enough to close the valves. This is commonly referred to as 'valve float', and it can result in piston to valve contact, severely damaging the engine. At high speeds the lubrication of piston cylinder wall interface tends to break down. This limits the piston speed for industrial engines to about 10 m/s.
Intake/exhaust port flow
The output power of an engine is dependent on the ability of intake (air–fuel mixture) and exhaust matter to move quickly through valve ports, typically located in the cylinder head. To increase an engine's output power, irregularities in the intake and exhaust paths, such as casting flaws, can be removed, and, with the aid of an air flow bench, the radii of valve port turns and valve seat configuration can be modified to reduce resistance. This process is called porting, and it can be done by hand or with a CNC machine.
Waste heat recovery of an internal combustion engine
An internal combustion engine is on average capable of converting only 40-45% of supplied energy into mechanical work. A large part of the waste energy is in the form of heat that is released to the environment through coolant, fins etc. If somehow waste heat could be captured and turned to mechanical energy, the engine's performance and/or fuel efficiency could be improved by improving the overall efficiency of the cycle. It has been found that even if 6% of the entirely wasted heat is recovered it can increase the engine efficiency greatly.
Many methods have been devised in order to extract waste heat out of an engine exhaust and use it further to extract some useful work, decreasing the exhaust pollutants at the same time. Use of the Rankine Cycle, turbocharging and thermoelectric generation can be very useful as a waste heat recovery system.
Supercharging
One way to increase engine power is to force more air into the cylinder so that more power can be produced from each power stroke. This can be done using some type of air compression device known as a supercharger, which can be powered by the engine crankshaft.
Supercharging increases the power output limits of an internal combustion engine relative to its displacement. Most commonly, the supercharger is always running, but there have been designs that allow it to be cut out or run at varying speeds (relative to engine speed). Mechanically driven supercharging has the disadvantage that some of the output power is used to drive the supercharger, while power is wasted in the high pressure exhaust, as the air has been compressed twice and then gains more potential volume in the combustion but it is only expanded in one stage.
Turbocharging
A turbocharger is a supercharger that is driven by the engine's exhaust gases, by means of a turbine. A turbocharger is incorporated into the exhaust system of a vehicle to make use of the expelled exhaust. It consists of a two piece, high-speed turbine assembly with one side that compresses the intake air, and the other side that is powered by the exhaust gas outflow.
When idling, and at low-to-moderate speeds, the turbine produces little power from the small exhaust volume, the turbocharger has little effect and the engine operates nearly in a naturally aspirated manner. When much more power output is required, the engine speed and throttle opening are increased until the exhaust gases are sufficient to 'spool up' the turbocharger's turbine to start compressing much more air than normal into the intake manifold. Thus, additional power (and speed) is expelled through the function of this turbine.
Turbocharging allows for more efficient engine operation because it is driven by exhaust pressure that would otherwise be (mostly) wasted, but there is a design limitation known as turbo lag. The increased engine power is not immediately available due to the need to sharply increase engine RPM, to build up pressure and to spin up the turbo, before the turbo starts to do any useful air compression. The increased intake volume causes increased exhaust and spins the turbo faster, and so forth until steady high power operation is reached. Another difficulty is that the higher exhaust pressure causes the exhaust gas to transfer more of its heat to the mechanical parts of the engine.
Rod and piston-to-stroke ratio
The rod-to-stroke ratio is the ratio of the length of the connecting rod to the length of the piston stroke. A longer rod reduces sidewise pressure of the piston on the cylinder wall and the stress forces, increasing engine life. It also increases the cost and engine height and weight.
A "square engine" is an engine with a bore diameter equal to its stroke length. An engine where the bore diameter is larger than its stroke length is an oversquare engine, conversely, an engine with a bore diameter that is smaller than its stroke length is an undersquare engine.
Valve train
The valves are typically operated by a camshaft rotating at half the speed of the crankshaft. It has a series of cams along its length, each designed to open a valve during the appropriate part of an intake or exhaust stroke. A tappet between valve and cam is a contact surface on which the cam slides to open the valve. Many engines use one or more camshafts "above" a row (or each row) of cylinders, as in the illustration, in which each cam directly actuates a valve through a flat tappet. In other engine designs the camshaft is in the crankcase, in which case each cam usually contacts a push rod, which contacts a rocker arm that opens a valve, or in case of a flathead engine a push rod is not necessary. The overhead cam design typically allows higher engine speeds because it provides the most direct path between cam and valve.
Valve clearance
Valve clearance refers to the small gap between a valve lifter and a valve stem that ensures that the valve completely closes. On engines with mechanical valve adjustment, excessive clearance causes noise from the valve train. A too-small valve clearance can result in the valves not closing properly. This results in a loss of performance and possibly overheating of exhaust valves. Typically, the clearance must be readjusted each with a feeler gauge.
Most modern production engines use hydraulic lifters to automatically compensate for valve train component wear. Dirty engine oil may cause lifter failure.
Energy balance
Otto engines are about 30% efficient; in other words, 30% of the energy generated by combustion is converted into useful rotational energy at the output shaft of the engine, while the remainder being lost due to waste heat, friction and engine accessories. There are a number of ways to recover some of the energy lost to waste heat. The use of a turbocharger in diesel engines is very effective by boosting incoming air pressure and in effect, provides the same increase in performance as having more displacement. The Mack Truck company, decades ago, developed a turbine system that converted waste heat into kinetic energy that it fed back into the engine's transmission. In 2005, BMW announced the development of the turbosteamer, a two-stage heat-recovery system similar to the Mack system that recovers 80% of the energy in the exhaust gas and raises the efficiency of an Otto engine by 15%. By contrast, a six-stroke engine may reduce fuel consumption by as much as 40%.
Modern engines are often intentionally built to be slightly less efficient than they could otherwise be. This is necessary for emission controls such as exhaust gas recirculation and catalytic converters that reduce smog and other atmospheric pollutants. Reductions in efficiency may be counteracted with an engine control unit using lean burn techniques.
In the United States, the Corporate Average Fuel Economy mandates that vehicles must achieve an average of compared to the current standard of . As automakers look to meet these standards by 2016, new ways of engineering the traditional internal combustion engine (ICE) have to be considered. Some potential solutions to increase fuel efficiency to meet new mandates include firing after the piston is farthest from the crankshaft, known as top dead centre, and applying the Miller cycle. Together, this redesign could significantly reduce fuel consumption and emissions.
Starting position, intake stroke, and compression stroke.
Ignition of fuel, power stroke, and exhaust stroke.
| Technology | Basics_8 | null |
170089 | https://en.wikipedia.org/wiki/Numerical%20integration | Numerical integration | In analysis, numerical integration comprises a broad family of algorithms for calculating the numerical value of a definite integral.
The term numerical quadrature (often abbreviated to quadrature) is more or less a synonym for "numerical integration", especially as applied to one-dimensional integrals. Some authors refer to numerical integration over more than one dimension as cubature; others take "quadrature" to include higher-dimensional integration.
The basic problem in numerical integration is to compute an approximate solution to a definite integral
to a given degree of accuracy. If is a smooth function integrated over a small number of dimensions, and the domain of integration is bounded, there are many methods for approximating the integral to the desired precision.
Numerical integration has roots in the geometrical problem of finding a square with the same area as a given plane figure (quadrature or squaring), as in the quadrature of the circle.
The term is also sometimes used to describe the numerical solution of differential equations.
Motivation and need
There are several reasons for carrying out numerical integration, as opposed to analytical integration by finding the antiderivative:
The integrand may be known only at certain points, such as obtained by sampling. Some embedded systems and other computer applications may need numerical integration for this reason.
A formula for the integrand may be known, but it may be difficult or impossible to find an antiderivative that is an elementary function. An example of such an integrand is , the antiderivative of which (the error function, times a constant) cannot be written in elementary form.
It may be possible to find an antiderivative symbolically, but it may be easier to compute a numerical approximation than to compute the antiderivative. That may be the case if the antiderivative is given as an infinite series or product, or if its evaluation requires a special function that is not available.
History
The term "numerical integration" first appears in 1915 in the publication A Course in Interpolation and Numeric Integration for the Mathematical Laboratory by David Gibb.
"Quadrature" is a historical mathematical term that means calculating area. Quadrature problems have served as one of the main sources of mathematical analysis. Mathematicians of Ancient Greece, according to the Pythagorean doctrine, understood calculation of area as the process of constructing geometrically a square having the same area (squaring). That is why the process was named "quadrature". For example, a quadrature of the circle, Lune of Hippocrates, The Quadrature of the Parabola. This construction must be performed only by means of compass and straightedge.
The ancient Babylonians used the trapezoidal rule to integrate the motion of Jupiter along the ecliptic.
For a quadrature of a rectangle with the sides a and b it is necessary to construct a square with the side (the Geometric mean of a and b). For this purpose it is possible to use the following fact: if we draw the circle with the sum of a and b as the diameter, then the height BH (from a point of their connection to crossing with a circle) equals their geometric mean. The similar geometrical construction solves a problem of a quadrature for a parallelogram and a triangle.
Problems of quadrature for curvilinear figures are much more difficult. The quadrature of the circle with compass and straightedge had been proved in the 19th century to be impossible. Nevertheless, for some figures (for example the Lune of Hippocrates) a quadrature can be performed. The quadratures of a sphere surface and a parabola segment done by Archimedes became the highest achievement of the antique analysis.
The area of the surface of a sphere is equal to quadruple the area of a great circle of this sphere.
The area of a segment of the parabola cut from it by a straight line is 4/3 the area of the triangle inscribed in this segment.
For the proof of the results Archimedes used the Method of exhaustion of Eudoxus.
In medieval Europe the quadrature meant calculation of area by any method. More often the Method of indivisibles was used; it was less rigorous, but more simple and powerful. With its help Galileo Galilei and Gilles de Roberval found the area of a cycloid arch, Grégoire de Saint-Vincent investigated the area under a hyperbola (Opus Geometricum, 1647), and Alphonse Antonio de Sarasa, de Saint-Vincent's pupil and commentator, noted the relation of this area to logarithms.
John Wallis algebrised this method: he wrote in his Arithmetica Infinitorum (1656) series that we now call the definite integral, and he calculated their values. Isaac Barrow and James Gregory made further progress: quadratures for some algebraic curves and spirals. Christiaan Huygens successfully performed a quadrature of some Solids of revolution.
The quadrature of the hyperbola by Saint-Vincent and de Sarasa provided a new function, the natural logarithm, of critical importance.
With the invention of integral calculus came a universal method for area calculation. In response, the term "quadrature" has become traditional, and instead the modern phrase "computation of a univariate definite integral" is more common.
Methods for one-dimensional integrals
A quadrature rule is an approximation of the definite integral of a function, usually stated as a weighted sum of function values at specified points within the domain of integration.
Numerical integration methods can generally be described as combining evaluations of the integrand to get an approximation to the integral. The integrand is evaluated at a finite set of points called integration points and a weighted sum of these values is used to approximate the integral. The integration points and weights depend on the specific method used and the accuracy required from the approximation.
An important part of the analysis of any numerical integration method is to study the behavior of the approximation error as a function of the number of integrand evaluations. A method that yields a small error for a small number of evaluations is usually considered superior. Reducing the number of evaluations of the integrand reduces the number of arithmetic operations involved, and therefore reduces the total error. Also, each evaluation takes time, and the integrand may be arbitrarily complicated.
Quadrature rules based on step functions
A "brute force" kind of numerical integration can be done, if the integrand is reasonably well-behaved (i.e. piecewise continuous and of bounded variation), by evaluating the integrand with very small increments.
This simplest method approximates the function by a step function (a piecewise constant function, or a segmented polynomial of degree zero) that passes through the point . This is called the midpoint rule or rectangle rule
Quadrature rules based on interpolating functions
A large class of quadrature rules can be derived by constructing interpolating functions that are easy to integrate. Typically these interpolating functions are polynomials. In practice, since polynomials of very high degree tend to oscillate wildly, only polynomials of low degree are used, typically linear and quadratic.
The interpolating function may be a straight line (an affine function, i.e. a polynomial of degree 1)
passing through the points and .
This is called the trapezoidal rule
For either one of these rules, we can make a more accurate approximation by breaking up the interval into some number of subintervals, computing an approximation for each subinterval, then adding up all the results. This is called a composite rule, extended rule, or iterated rule. For example, the composite trapezoidal rule can be stated as
where the subintervals have the form with and Here we used subintervals of the same length but one could also use intervals of varying length .
Interpolation with polynomials evaluated at equally spaced points in yields the Newton–Cotes formulas, of which the rectangle rule and the trapezoidal rule are examples. Simpson's rule, which is based on a polynomial of order 2, is also a Newton–Cotes formula.
Quadrature rules with equally spaced points have the very convenient property of nesting. The corresponding rule with each interval subdivided includes all the current points, so those integrand values can be re-used.
If we allow the intervals between interpolation points to vary, we find another group of quadrature formulas, such as the Gaussian quadrature formulas. A Gaussian quadrature rule is typically more accurate than a Newton–Cotes rule that uses the same number of function evaluations, if the integrand is smooth (i.e., if it is sufficiently differentiable). Other quadrature methods with varying intervals include Clenshaw–Curtis quadrature (also called Fejér quadrature) methods, which do nest.
Gaussian quadrature rules do not nest, but the related Gauss–Kronrod quadrature formulas do.
Adaptive algorithms
Extrapolation methods
The accuracy of a quadrature rule of the Newton–Cotes type is generally a function of the number of evaluation points. The result is usually more accurate as the number of evaluation points increases, or, equivalently, as the width of the step size between the points decreases. It is natural to ask what the result would be if the step size were allowed to approach zero. This can be answered by extrapolating the result from two or more nonzero step sizes, using series acceleration methods such as Richardson extrapolation. The extrapolation function may be a polynomial or rational function. Extrapolation methods are described in more detail by Stoer and Bulirsch (Section 3.4) and are implemented in many of the routines in the QUADPACK library.
Conservative (a priori) error estimation
Let have a bounded first derivative over i.e. The mean value theorem for where gives
for some depending on .
If we integrate in from to on both sides and take the absolute values, we obtain
We can further approximate the integral on the right-hand side by bringing the absolute value into the integrand, and replacing the term in by an upper bound
where the supremum was used to approximate.
Hence, if we approximate the integral by the quadrature rule our error is no greater than the right hand side of . We can convert this into an error analysis for the Riemann sum, giving an upper bound of
for the error term of that particular approximation. (Note that this is precisely the error we calculated for the example .) Using more derivatives, and by tweaking the quadrature, we can do a similar error analysis using a Taylor series (using a partial sum with remainder term) for f. This error analysis gives a strict upper bound on the error, if the derivatives of f are available.
This integration method can be combined with interval arithmetic to produce computer proofs and verified calculations.
Integrals over infinite intervals
Several methods exist for approximate integration over unbounded intervals. The standard technique involves specially derived quadrature rules, such as Gauss-Hermite quadrature for integrals on the whole real line and Gauss-Laguerre quadrature for integrals on the positive reals. Monte Carlo methods can also be used, or a change of variables to a finite interval; e.g., for the whole line one could use
and for semi-infinite intervals one could use
as possible transformations.
Multidimensional integrals
The quadrature rules discussed so far are all designed to compute one-dimensional integrals. To compute integrals in multiple dimensions, one approach is to phrase the multiple integral as repeated one-dimensional integrals by applying Fubini's theorem (the tensor product rule). This approach requires the function evaluations to grow exponentially as the number of dimensions increases. Three methods are known to overcome this so-called curse of dimensionality.
A great many additional techniques for forming multidimensional cubature integration rules for a variety of weighting functions are given in the monograph by Stroud.
Integration on the sphere has been reviewed by Hesse et al. (2015).
Monte Carlo
Monte Carlo methods and quasi-Monte Carlo methods are easy to apply to multi-dimensional integrals. They may yield greater accuracy for the same number of function evaluations than repeated integrations using one-dimensional methods.
A large class of useful Monte Carlo methods are the so-called Markov chain Monte Carlo algorithms, which include the Metropolis–Hastings algorithm and Gibbs sampling.
Sparse grids
Sparse grids were originally developed by Smolyak for the quadrature of high-dimensional functions. The method is always based on a one-dimensional quadrature rule, but performs a more sophisticated combination of univariate results. However, whereas the tensor product rule guarantees that the weights of all of the cubature points will be positive if the weights of the quadrature points were positive, Smolyak's rule does not guarantee that the weights will all be positive.
Bayesian quadrature
Bayesian quadrature is a statistical approach to the numerical problem of computing integrals and falls under the field of probabilistic numerics. It can provide a full handling of the uncertainty over the solution of the integral expressed as a Gaussian process posterior variance.
Connection with differential equations
The problem of evaluating the definite integral
can be reduced to an initial value problem for an ordinary differential equation by applying the first part of the fundamental theorem of calculus. By differentiating both sides of the above with respect to the argument x, it is seen that the function F satisfies
Numerical methods for ordinary differential equations, such as Runge–Kutta methods, can be applied to the restated problem and thus be used to evaluate the integral. For instance, the standard fourth-order Runge–Kutta method applied to the differential equation yields Simpson's rule from above.
The differential equation has a special form: the right-hand side contains only the independent variable (here ) and not the dependent variable (here ). This simplifies the theory and algorithms considerably. The problem of evaluating integrals is thus best studied in its own right.
Conversely, the term "quadrature" may also be used for the solution of differential equations: "solving by quadrature" or "reduction to quadrature" means expressing its solution in terms of integrals.
| Mathematics | Discrete mathematics | null |
170097 | https://en.wikipedia.org/wiki/Mean%20free%20path | Mean free path | In physics, mean free path is the average distance over which a moving particle (such as an atom, a molecule, or a photon) travels before substantially changing its direction or energy (or, in a specific context, other properties), typically as a result of one or more successive collisions with other particles.
Scattering theory
Imagine a beam of particles being shot through a target, and consider an infinitesimally thin slab of the target (see the figure). The atoms (or particles) that might stop a beam particle are shown in red. The magnitude of the mean free path depends on the characteristics of the system. Assuming that all the target particles are at rest but only the beam particle is moving, that gives an expression for the mean free path:
where is the mean free path, is the number of target particles per unit volume, and is the effective cross-sectional area for collision.
The area of the slab is , and its volume is . The typical number of stopping atoms in the slab is the concentration times the volume, i.e., . The probability that a beam particle will be stopped in that slab is the net area of the stopping atoms divided by the total area of the slab:
where is the area (or, more formally, the "scattering cross-section") of one atom.
The drop in beam intensity equals the incoming beam intensity multiplied by the probability of the particle being stopped within the slab:
This is an ordinary differential equation:
whose solution is known as Beer–Lambert law and has the form , where is the distance traveled by the beam through the target, and is the beam intensity before it entered the target; is called the mean free path because it equals the mean distance traveled by a beam particle before being stopped. To see this, note that the probability that a particle is absorbed between and is given by
Thus the expectation value (or average, or simply mean) of is
The fraction of particles that are not stopped (attenuated) by the slab is called transmission , where is equal to the thickness of the slab.
Kinetic theory of gases
In the kinetic theory of gases, the mean free path of a particle, such as a molecule, is the average distance the particle travels between collisions with other moving particles. The derivation above assumed the target particles to be at rest; therefore, in reality, the formula holds for a beam particle with a high speed relative to the velocities of an ensemble of identical particles with random locations. In that case, the motions of target particles are comparatively negligible, hence the relative velocity .
If, on the other hand, the beam particle is part of an established equilibrium with identical particles, then the square of relative velocity is:
In equilibrium, and are random and uncorrelated, therefore , and the relative speed is
This means that the number of collisions is times the number with stationary targets. Therefore, the following relationship applies:
and using (ideal gas law) and (effective cross-sectional area for spherical particles with diameter ), it may be shown that the mean free path is
where k is the Boltzmann constant, is the pressure of the gas and is the absolute temperature.
In practice, the diameter of gas molecules is not well defined. In fact, the kinetic diameter of a molecule is defined in terms of the mean free path. Typically, gas molecules do not behave like hard spheres, but rather attract each other at larger distances and repel each other at shorter distances, as can be described with a Lennard-Jones potential. One way to deal with such "soft" molecules is to use the Lennard-Jones σ parameter as the diameter.
Another way is to assume a hard-sphere gas that has the same viscosity as the actual gas being considered. This leads to a mean free path
where is the molecular mass, is the density of ideal gas, and μ is the dynamic viscosity. This expression can be put into the following convenient form
with being the specific gas constant, equal to 287 J/(kg*K) for air.
The following table lists some typical values for air at different pressures at room temperature. Note that different definitions of the molecular diameter, as well as different assumptions about the value of atmospheric pressure (100 vs 101.3 kPa) and room temperature (293.17 K vs 296.15 K or even 300 K) can lead to slightly different values of the mean free path.
In other fields
Radiography
In gamma-ray radiography the mean free path of a pencil beam of mono-energetic photons is the average distance a photon travels between collisions with atoms of the target material. It depends on the material and the energy of the photons:
where μ is the linear attenuation coefficient, μ/ρ is the mass attenuation coefficient and ρ is the density of the material. The mass attenuation coefficient can be looked up or calculated for any material and energy combination using the National Institute of Standards and Technology (NIST) databases.
In X-ray radiography the calculation of the mean free path is more complicated, because photons are not mono-energetic, but have some distribution of energies called a spectrum. As photons move through the target material, they are attenuated with probabilities depending on their energy, as a result their distribution changes in process called spectrum hardening. Because of spectrum hardening, the mean free path of the X-ray spectrum changes with distance.
Sometimes one measures the thickness of a material in the number of mean free paths. Material with the thickness of one mean free path will attenuate to 37% (1/e) of photons. This concept is closely related to half-value layer (HVL): a material with a thickness of one HVL will attenuate 50% of photons. A standard x-ray image is a transmission image, an image with negative logarithm of its intensities is sometimes called a number of mean free paths image.
Electronics
In macroscopic charge transport, the mean free path of a charge carrier in a metal is proportional to the electrical mobility , a value directly related to electrical conductivity, that is:
where q is the charge, is the mean free time, m* is the effective mass, and vF is the Fermi velocity of the charge carrier. The Fermi velocity can easily be derived from the Fermi energy via the non-relativistic kinetic energy equation. In thin films, however, the film thickness can be smaller than the predicted mean free path, making surface scattering much more noticeable, effectively increasing the resistivity.
Electron mobility through a medium with dimensions smaller than the mean free path of electrons occurs through ballistic conduction or ballistic transport. In such scenarios electrons alter their motion only in collisions with conductor walls.
Optics
If one takes a suspension of non-light-absorbing particles of diameter d with a volume fraction Φ, the mean free path of the photons is:
where Qs is the scattering efficiency factor. Qs can be evaluated numerically for spherical particles using Mie theory.
Acoustics
In an otherwise empty cavity, the mean free path of a single particle bouncing off the walls is:
where V is the volume of the cavity, S is the total inside surface area of the cavity, and F is a constant related to the shape of the cavity. For most simple cavity shapes, F is approximately 4.
This relation is used in the derivation of the Sabine equation in acoustics, using a geometrical approximation of sound propagation.
Nuclear and particle physics
In particle physics the concept of the mean free path is not commonly used, being replaced by the similar concept of attenuation length. In particular, for high-energy photons, which mostly interact by electron–positron pair production, the radiation length is used much like the mean free path in radiography.
Independent-particle models in nuclear physics require the undisturbed orbiting of nucleons within the nucleus before they interact with other nucleons.
| Physical sciences | Thermodynamics | Physics |
170104 | https://en.wikipedia.org/wiki/Juniper | Juniper | Junipers are coniferous trees and shrubs in the genus Juniperus ( ) of the cypress family Cupressaceae. Depending on the taxonomy, between 50 and 67 species of junipers are widely distributed throughout the Northern Hemisphere as far south as tropical Africa, including the Arctic, parts of Asia, and Central America. The highest-known juniper forest occurs at an altitude of in southeastern Tibet and the northern Himalayas, creating one of the highest tree lines on earth.
Description
Junipers vary in size and shape from tall trees, tall, to columnar or low-spreading shrubs with long, trailing branches. They are evergreen with needle-like and/or scale-like leaves. They can be either monoecious or dioecious. The female seed cones are very distinctive, with fleshy, fruit-like coalescing scales which fuse together to form a berrylike structure (galbulus), long, with one to 12 unwinged, hard-shelled seeds. In some species, these "berries" are red-brown or orange, but in most, they are blue; they are often aromatic and can be used as a spice. The seed maturation time varies between species from 6 to 18 months after pollination. The male cones are similar to the other Cupressaceae, with 6 to 20 scales.
In hardiness zones 7 through 10, junipers can bloom and release pollen several times each year. Different junipers bloom in autumn, while most pollinate from early winter until late spring.
Many junipers (e.g. J. chinensis, J. virginiana) have two types of leaves; seedlings and some twigs of older trees have needle-like leaves long, on mature plants the leaves are overlapping like (mostly) tiny scales, measuring . When juvenile foliage occurs on mature plants, it is most often found on shaded shoots, with adult foliage in full sunlight. Leaves on fast-growing 'whip' shoots are often intermediate between juvenile and adult.
In some species (e.g. J. communis, J. squamata), all the foliage is of the juvenile needle-like type, with no scale leaves. In some of these (e.g. J. communis), the needles are jointed at the base, while in others (e.g. J. squamata), the needles merge smoothly with the stem.
The needle leaves of junipers are hard and sharp, making the juvenile foliage very prickly to handle. This can be a valuable identification feature in seedlings, as the otherwise very similar juvenile foliage of cypresses (Cupressus, Chamaecyparis) and other related genera are soft and not prickly.
Junipers are gymnosperms, which means they have seeds, but no flowers or fruits. Depending on the species, the seeds they produce take 1–3 years to develop. The impermeable coat of the seed keeps water from getting in and protects the embryo when dispersed. It can also result in a long dormancy that is usually broken by physically damaging the seed coat. Dispersal can occur from being swallowed whole by frugivores and mammals. The resistance of the seed coat allows it to be passed down through the digestive system without being destroyed along the way. These seeds last a long time, as they can be dispersed long distances over the course of a few years.
Classification
Sections
The genus has been divided into sections in somewhat different ways. A system based on molecular phylogenetic data from 2013 and earlier used three sections:
Section Caryocedrus – 1 species with large, blue, woody, 3-seeded cones; native to the Mediterranean
Section Juniperus – 14 species with blue or red seed cones, often with 3 seeds; 12 species native to the eastern hemisphere, one endemic to North America, and one species, J. communis, circumboreal
Section Sabina – about 60 species with variously coloured seed cones with 1 to 13 seeds; species about equally divided between the eastern and western hemispheres Juniperus sect. Sabina was further divided into clades.
A new classification of gymnosperms published in 2022 recognised the sections as three separate genera: Arceuthos for section Caryocedrus, Sabina for section Sabina, and Juniperus sensu stricto for section Juniperus.
Species
Juniperus sect. Caryocedrus
Cones with three seeds fused together; needles with two stomatal bands.
One species:
Juniperus drupacea – Syrian juniper
Juniperus sect. Juniperus
Needle-leaf junipers; the adult leaves are needle-like, in whorls of three, and jointed at the base. Species:
Juniperus sect. Juniperus subsect. Juniperus – cones with three separate seeds; needles with one stomatal band
Juniperus communis – common juniper
Juniperus communis subsp. alpina – alpine juniper
Juniperus conferta, syn. Juniperus rigida var. conferta (Parl.) Patschke – shore juniper
Juniperus rigida – Temple juniper or needle juniper
Juniperus sect. Juniperus subsect. Oxycedrus – cones with three separate seeds; needles with two stomatal bands
Juniperus brevifolia – Azores juniper
Juniperus cedrus – Canary Islands juniper
Juniperus formosana – Chinese prickly juniper
Juniperus lutchuensis, syn. Juniperus taxifolia var. lutchuensis (Koidz.) Satake – Ryukyu juniper
Juniperus oxycedrus – Western prickly juniper, cade juniper
Juniperus macrocarpa – large-berry juniper
Juniperus sect. Sabina
Scale-leaf junipers; adult leaves are mostly scale-like, similar to those of Cupressus species, in opposite pairs or whorls of three, and the juvenile needle-like leaves are not jointed at the base (including in the few that have only needle-like leaves; see below right).
Old World species
Juniperus chinensis – Chinese juniper
Juniperus convallium – Mekong juniper
Juniperus excelsa – Greek juniper
Juniperus foetidissima – stinking juniper
Juniperus indica – Himalayan black juniper
Juniperus komarovii – Komarov's juniper
Juniperus phoenicea – Phoenicean juniper
Juniperus pingii – Ping juniper
Juniperus procera – East African juniper
Juniperus procumbens – Ibuki juniper
Juniperus pseudosabina – Xinjiang juniper
Juniperus recurva – Himalayan juniper
Juniperus sabina – Savin juniper
Juniperus saltuaria – Sichuan juniper
Juniperus semiglobosa – Himalayan pencil juniper
Juniperus seravschanica – Pashtun juniper
Juniperus squamata – flaky juniper
Juniperus thurifera – Spanish juniper
Juniperus tibetica – Tibetan juniper
New World species
Juniperus angosturana – Mexican one-seed juniper
Juniperus ashei – Ashe juniper
Juniperus arizonica – redberry juniper, roseberry juniper
Juniperus barbadensis – West Indies juniper
Juniperus bermudiana – Bermuda juniper
Juniperus blancoi – Blanco's juniper
Juniperus californica – California juniper
Juniperus coahuilensis – Coahuila juniper
Juniperus comitana – Comitán juniper
Juniperus deppeana – alligator juniper
Juniperus durangensis – Durango juniper
Juniperus flaccida – Mexican weeping juniper
Juniperus gamboana – Gamboa juniper
Juniperus grandis – Sierra juniper
Juniperus horizontalis – creeping juniper
Juniperus jaliscana – Jalisco juniper
Juniperus maritima, syn. Juniperus scopulorum – seaside juniper
Juniperus monosperma – one-seed juniper
Juniperus monticola – mountain juniper
Juniperus occidentalis – western juniper
Juniperus osteosperma – Utah juniper
Juniperus pinchotii – Pinchot juniper
Juniperus saltillensis – Saltillo juniper
Juniperus scopulorum – Rocky Mountain juniper
Juniperus standleyi – Standley's juniper
Juniperus virginiana – eastern juniper, eastern redcedar
Juniperus virginiana subsp. silicicola – Southern juniper
Juniperus zanonii (proposed)
Additional species
, Plants of the World Online accepts the following additional species to those listed above:
Juniperus canariensis Guyot & Mathou
Juniperus coxii A.B.Jacks.
Juniperus deltoides R.P.Adams
Juniperus gracilior Pilg.
Juniperus mairei Lemée & H.Lév.
Juniperus morrisonicola Hayata
Juniperus mucronata R.P.Adams
Juniperus navicularis Gand.
Juniperus poblana (Martínez) R.P.Adams
Juniperus polycarpos K.Koch
Juniperus przewalskii Kom.
Juniperus saxicola Britton & P.Wilson
Juniperus taxifolia Hook. & Arn.
Juniperus tsukusiensis Masam.
Juniperus turbinata Guss.
Ecology
Juniper plants thrive in a variety of environments. The junipers from Lahaul valley can be found in dry, rocky locations planted in stony soils. Grazing animals and the villagers are rapidly using up these plants. There are several important features of the leaves and wood of this plant that cause villagers to cut down these trees and make use of them. Additionally, the western juniper plants, a particular species in the juniper genus, are found in woodlands where there are large, open spaces. Junipers are known to encompass open areas so that they have more exposure to rainfall. Decreases in fires and a lack of livestock grazing are the two major causes of western juniper takeover. This invasion of junipers is driving changes in the environment. For instance, the ecosystem for other species previously living in the environment and farm animals has been compromised. When junipers increase in population, there is a decrease in woody species like mountain big sagebrush and aspen. Among the juniper trees themselves, there is increased competition, which results in a decrease in berry production. Herbaceous cover decreases, and junipers are often mistaken for weeds. As a result, several farmers have thinned the juniper trees or removed them completely. However, this reduction did not result in any significant difference on wildlife survival. Some small mammals found it advantageous to have thinner juniper trees, while cutting down the entire tree was not favorable.
Some junipers are susceptible to Gymnosporangium rust disease and can be a serious problem for those people growing apple trees, an alternate host of the disease.
Juniper is the exclusive food plant of the larvae of some moths and butterflies, including Bucculatrix inusitata, juniper carpet, Chionodes electella, Chionodes viduella, juniper pug, and pine beauty. Those of the tortrix moth Cydia duplicana feed on the bark around injuries or canker.
Cultivation
Junipers are among the most popular conifers to be cultivated as ornamental subjects for parks and gardens. They have been bred over many years to produce a wide range of forms, in terms of colour, shape and size. They include some of the dwarfest (miniature) cultivars. They are also used for bonsai. Some species found in cultivation include:
Juniperus chinensis
Juniperus communis
Juniperus horizontalis
Juniperus × pfitzeriana
Juniperus procumbens
Juniperus rigida
Juniperus scopulorum
Juniperus squamata
Toxicity
In drier areas, juniper pollen easily becomes airborne and can be inhaled into the lungs. This pollen can also irritate the skin and cause contact dermatitis. Cross-allergenic reactions are common between juniper pollen and the pollen of all species of cypress.
Monoecious juniper plants are highly allergenic, with an Ogren Plant Allergy Scale (OPALS) rating of 9 out of 10. Completely male juniper plants have an OPALS rating of 10, and release abundant amounts of pollen. Conversely, all-female juniper plants have an OPALS rating of 1, and are considered "allergy-fighting".
Uses
Ethnic and herbal use
Most species of juniper are flexible and have a high compression strength-to-weight ratio. This has made the wood a traditional choice for the construction of hunting bows among some of the Native American cultures in the Great Basin region. These bow staves are typically backed with sinew to provide tension strength that the wood may lack.
Ancient Mesopotamians believed that juniper oil could be used to ward off the evil eye.
Embalming vessels in the burial chambers from a 26th Dynasty embalming workshop at Saqqara have shown the usage of Juniper oil/tar.
Some Indigenous peoples of the Americas use juniper in traditional medicine; for instance the Dineh (Navajo), who use it for diabetes. Juniper ash has also been historically consumed as a source of calcium by the Navajo people.
Juniper is traditionally used in Scottish folkloric and Gaelic Polytheist saining rites, such as those performed at Hogmanay (New Year), where the smoke of burning juniper, accompanied by traditional prayers and other customary rites, is used to cleanse, bless, and protect the household and its inhabitants.
Local people in Lahaul Valley present juniper leaves to their deities as a folk tradition. It is also useful as a folk remedy for pains and aches, as well as epilepsy and asthma. They are reported to collect large amounts of juniper leaves and wood for building and religious purposes.
General use
Juniper berries are a spice used in a wide variety of culinary dishes and are best known for the primary flavoring in gin (and responsible for gin's name, which is a shortening of the Dutch word for juniper, jenever). A juniper-based spirit is made by fermenting juniper berries and water to create a "wine" that is then distilled. This is often sold as a juniper brandy in eastern Europe. Juniper berries are also used as the primary flavor in the liquor jenever. Juniper berry sauce is often a popular flavoring choice for quail, pheasant, veal, rabbit, venison, and other game dishes.
A tea can be made from the young twigs. Twigs or needles are used to flavour the traditional Finnish junperbeer, sahti as well.
Dense and rot resistant, the irregular trunks of junipers have been used as fence posts and firewood. Stands that produce enough wood for specialty uses generally go under the common name "cedar", including Juniperus virginiana, the "red cedar" that is used widely in cedar drawers and closets. The lack of space or a hyphen between the words "red" and "cedar" is sometimes used to indicate that this species is not a true cedar (Cedrus).
Juniper in weave is a traditional cladding technique used in Northern Europe, e.g. at Havrå, Norway.
Juniper berries are steam distilled to produce an essential oil that may vary from colorless to yellow or pale green. Some of its chemical components are terpenoids and aromatic compounds, such as cadinene, a sesquiterpene.
| Biology and health sciences | Gymnosperms | null |
170165 | https://en.wikipedia.org/wiki/Fermi%20energy | Fermi energy | The Fermi energy is a concept in quantum mechanics usually referring to the energy difference between the highest and lowest occupied single-particle states in a quantum system of non-interacting fermions at absolute zero temperature.
In a Fermi gas, the lowest occupied state is taken to have zero kinetic energy, whereas in a metal, the lowest occupied state is typically taken to mean the bottom of the conduction band.
The term "Fermi energy" is often used to refer to a different yet closely related concept, the Fermi level (also called electrochemical potential).
There are a few key differences between the Fermi level and Fermi energy, at least as they are used in this article:
The Fermi energy is only defined at absolute zero, while the Fermi level is defined for any temperature.
The Fermi energy is an energy difference (usually corresponding to a kinetic energy), whereas the Fermi level is a total energy level including kinetic energy and potential energy.
The Fermi energy can only be defined for non-interacting fermions (where the potential energy or band edge is a static, well defined quantity), whereas the Fermi level remains well defined even in complex interacting systems, at thermodynamic equilibrium.
Since the Fermi level in a metal at absolute zero is the energy of the highest occupied single particle state,
then the Fermi energy in a metal is the energy difference between the Fermi level and lowest occupied single-particle state, at zero-temperature.
Context
In quantum mechanics, a group of particles known as fermions (for example, electrons, protons and neutrons) obey the Pauli exclusion principle. This states that two fermions cannot occupy the same quantum state. Since an idealized non-interacting Fermi gas can be analyzed in terms of single-particle stationary states, we can thus say that two fermions cannot occupy the same stationary state. These stationary states will typically be distinct in energy. To find the ground state of the whole system, we start with an empty system, and add particles one at a time, consecutively filling up the unoccupied stationary states with the lowest energy. When all the particles have been put in, the Fermi energy is the kinetic energy of the highest occupied state.
As a consequence, even if we have extracted all possible energy from a Fermi gas by cooling it to near absolute zero temperature, the fermions are still moving around at a high speed. The fastest ones are moving at a velocity corresponding to a kinetic energy equal to the Fermi energy. This speed is known as the Fermi velocity. Only when the temperature exceeds the related Fermi temperature, do the particles begin to move significantly faster than at absolute zero.
The Fermi energy is an important concept in the solid state physics of metals and superconductors. It is also a very important quantity in the physics of quantum liquids like low temperature helium (both normal and superfluid 3He), and it is quite important to nuclear physics and to understanding the stability of white dwarf stars against gravitational collapse.
Formula and typical values
The Fermi energy for a three-dimensional, non-relativistic, non-interacting ensemble of identical spin- fermions is given by
where N is the number of particles, m0 the rest mass of each fermion, V the volume of the system, and the reduced Planck constant.
Metals
Under the free electron model, the electrons in a metal can be considered to form a Fermi gas. The number density of conduction electrons in metals ranges between approximately 1028 and 1029 electrons/m3, which is also the typical density of atoms in ordinary solid matter. This number density produces a Fermi energy of the order of 2 to 10 electronvolts.
White dwarfs
Stars known as white dwarfs have mass comparable to the Sun, but have about a hundredth of its radius. The high densities mean that the electrons are no longer bound to single nuclei and instead form a degenerate electron gas. Their Fermi energy is about 0.3 MeV.
Nucleus
Another typical example is that of the nucleons in the nucleus of an atom. The radius of the nucleus admits deviations, so a typical value for the Fermi energy is usually given as 38 MeV.
Related quantities
Using this definition of above for the Fermi energy, various related quantities can be useful.
The Fermi temperature is defined as
where is the Boltzmann constant, and the Fermi energy. The Fermi temperature can be thought of as the temperature at which thermal effects are comparable to quantum effects associated with Fermi statistics. The Fermi temperature for a metal is a couple of orders of magnitude above room temperature.
Other quantities defined in this context are Fermi momentum
and Fermi velocity
These quantities are respectively the momentum and group velocity of a fermion at the Fermi surface.
The Fermi momentum can also be described as
where , called the Fermi wavevector, is the radius of the Fermi sphere.
is the electron density.
These quantities may not be well-defined in cases where the Fermi surface is non-spherical.
| Physical sciences | Statistical mechanics | Physics |
170167 | https://en.wikipedia.org/wiki/Maxwell%E2%80%93Boltzmann%20statistics | Maxwell–Boltzmann statistics | In statistical mechanics, Maxwell–Boltzmann statistics describes the distribution of classical material particles over various energy states in thermal equilibrium. It is applicable when the temperature is high enough or the particle density is low enough to render quantum effects negligible.
The expected number of particles with energy for Maxwell–Boltzmann statistics is
where:
is the energy of the i-th energy level,
is the average number of particles in the set of states with energy ,
is the degeneracy of energy level i, that is, the number of states with energy which may nevertheless be distinguished from each other by some other means,
μ is the chemical potential,
k is the Boltzmann constant,
T is absolute temperature,
N is the total number of particles:
Z is the partition function:
e is Euler's number
Equivalently, the number of particles is sometimes expressed as
where the index i now specifies a particular state rather than the set of all states with energy , and .
History
Maxwell–Boltzmann statistics grew out of the Maxwell–Boltzmann distribution, most likely as a distillation of the underlying technique. The distribution was first derived by Maxwell in 1860 on heuristic grounds. Boltzmann later, in the 1870s, carried out significant investigations into the physical origins of this distribution. The distribution can be derived on the ground that it maximizes the entropy of the system.
Applicability
Maxwell–Boltzmann statistics is used to derive the Maxwell–Boltzmann distribution of an ideal gas. However, it can also be used to extend that distribution to particles with a different energy–momentum relation, such as relativistic particles (resulting in Maxwell–Jüttner distribution), and to other than three-dimensional spaces.
Maxwell–Boltzmann statistics is often described as the statistics of "distinguishable" classical particles. In other words, the configuration of particle A in state 1 and particle B in state 2 is different from the case in which particle B is in state 1 and particle A is in state 2. This assumption leads to the proper (Boltzmann) statistics of particles in the energy states, but yields non-physical results for the entropy, as embodied in the Gibbs paradox.
At the same time, there are no real particles that have the characteristics required by Maxwell–Boltzmann statistics. Indeed, the Gibbs paradox is resolved if we treat all particles of a certain type (e.g., electrons, protons,photon etc.) as principally indistinguishable. Once this assumption is made, the particle statistics change. The change in entropy in the entropy of mixing example may be viewed as an example of a non-extensive entropy resulting from the distinguishability of the two types of particles being mixed.
Quantum particles are either bosons (following Bose–Einstein statistics) or fermions (subject to the Pauli exclusion principle, following instead Fermi–Dirac statistics). Both of these quantum statistics approach the Maxwell–Boltzmann statistics in the limit of high temperature and low particle density.
Derivations
Maxwell–Boltzmann statistics can be derived in various statistical mechanical thermodynamic ensembles:
The grand canonical ensemble, exactly.
The canonical ensemble, exactly.
The microcanonical ensemble, but only in the thermodynamic limit.
In each case it is necessary to assume that the particles are non-interacting, and that multiple particles can occupy the same state and do so independently.
Derivation from microcanonical ensemble
Suppose we have a container with a huge number of very small particles all with identical physical characteristics (such as mass, charge, etc.). Let's refer to this as the system. Assume that though the particles have identical properties, they are distinguishable. For example, we might identify each particle by continually observing their trajectories, or by placing a marking on each one, e.g., drawing a different number on each one as is done with lottery balls.
The particles are moving inside that container in all directions with great speed. Because the particles are speeding around, they possess some energy. The Maxwell–Boltzmann distribution is a mathematical function that describes about how many particles in the container have a certain energy. More precisely, the Maxwell–Boltzmann distribution gives the non-normalized probability (this means that the probabilities do not add up to 1) that the state corresponding to a particular energy is occupied.
In general, there may be many particles with the same amount of energy . Let the number of particles with the same energy be , the number of particles possessing another energy be , and so forth for all the possible energies To describe this situation, we say that is the occupation number of the energy level If we know all the occupation numbers then we know the total energy of the system. However, because we can distinguish between which particles are occupying each energy level, the set of occupation numbers does not completely describe the state of the system. To completely describe the state of the system, or the microstate, we must specify exactly which particles are in each energy level. Thus when we count the number of possible states of the system, we must count each and every microstate, and not just the possible sets of occupation numbers.
To begin with, assume that there is only one state at each energy level (there is no degeneracy). What follows next is a bit of combinatorial thinking which has little to do in accurately describing the reservoir of particles. For instance, let's say there is a total of boxes labelled . With the concept of combination, we could calculate how many ways there are to arrange into the set of boxes, where the order of balls within each box isn’t tracked. First, we select balls from a total of balls to place into box , and continue to select for each box from the remaining balls, ensuring that every ball is placed in one of the boxes. The total number of ways that the balls can be arranged is
As every ball has been placed into a box, , and we simplify the expression as
This is just the multinomial coefficient, the number of ways of arranging N items into k boxes, the l-th box holding Nl items, ignoring the permutation of items in each box.
Now, consider the case where there is more than one way to put particles in the box (i.e. taking the degeneracy problem into consideration). If the -th box has a "degeneracy" of , that is, it has "sub-boxes" ( boxes with the same energy . These states/boxes with the same energy are called degenerate states.), such that any way of filling the -th box where the number in the sub-boxes is changed is a distinct way of filling the box, then the number of ways of filling the i-th box must be increased by the number of ways of distributing the objects in the "sub-boxes". The number of ways of placing distinguishable objects in "sub-boxes" is (the first object can go into any of the boxes, the second object can also go into any of the boxes, and so on). Thus the number of ways that a total of particles can be classified into energy levels according to their energies, while each level having distinct states such that the i-th level accommodates particles is:
This is the form for W first derived by Boltzmann. Boltzmann's fundamental equation relates the thermodynamic entropy S to the number of microstates W, where k is the Boltzmann constant. It was pointed out by Gibbs however, that the above expression for W does not yield an extensive entropy, and is therefore faulty. This problem is known as the Gibbs paradox. The problem is that the particles considered by the above equation are not indistinguishable. In other words, for two particles (A and B) in two energy sublevels the population represented by [A,B] is considered distinct from the population [B,A] while for indistinguishable particles, they are not. If we carry out the argument for indistinguishable particles, we are led to the Bose–Einstein expression for W:
The Maxwell–Boltzmann distribution follows from this Bose–Einstein distribution for temperatures well above absolute zero, implying that . The Maxwell–Boltzmann distribution also requires low density, implying that . Under these conditions, we may use Stirling's approximation for the factorial:
to write:
Using the fact that for we can again use Stirling's approximation to write:
This is essentially a division by N! of Boltzmann's original expression for W, and this correction is referred to as .
We wish to find the for which the function is maximized, while considering the constraint that there is a fixed number of particles and a fixed energy in the container. The maxima of and are achieved by the same values of and, since it is easier to accomplish mathematically, we will maximize the latter function instead. We constrain our solution using Lagrange multipliers forming the function:
Finally
In order to maximize the expression above we apply Fermat's theorem (stationary points), according to which local extrema, if exist, must be at critical points (partial derivatives vanish):
By solving the equations above () we arrive to an expression for :
Substituting this expression for into the equation for and assuming that yields:
or, rearranging:
Boltzmann realized that this is just an expression of the Euler-integrated fundamental equation of thermodynamics. Identifying E as the internal energy, the Euler-integrated fundamental equation states that :
where T is the temperature, P is pressure, V is volume, and μ is the chemical potential. Boltzmann's equation is the realization that the entropy is proportional to with the constant of proportionality being the Boltzmann constant. Using the ideal gas equation of state (PV = NkT), It follows immediately that and so that the populations may now be written:
Note that the above formula is sometimes written:
where is the absolute activity.
Alternatively, we may use the fact that
to obtain the population numbers as
where Z is the partition function defined by:
In an approximation where εi is considered to be a continuous variable, the Thomas–Fermi approximation yields a continuous degeneracy g proportional to so that:
which is just the Maxwell–Boltzmann distribution for the energy.
Derivation from canonical ensemble
In the above discussion, the Boltzmann distribution function was obtained via directly analysing the multiplicities of a system. Alternatively, one can make use of the canonical ensemble. In a canonical ensemble, a system is in thermal contact with a reservoir. While energy is free to flow between the system and the reservoir, the reservoir is thought to have infinitely large heat capacity as to maintain constant temperature, T, for the combined system.
In the present context, our system is assumed to have the energy levels with degeneracies . As before, we would like to calculate the probability that our system has energy .
If our system is in state , then there would be a corresponding number of microstates available to the reservoir. Call this number . By assumption, the combined system (of the system we are interested in and the reservoir) is isolated, so all microstates are equally probable. Therefore, for instance, if , we can conclude that our system is twice as likely to be in state than . In general, if is the probability that our system is in state ,
Since the entropy of the reservoir , the above becomes
Next we recall the thermodynamic identity (from the first law of thermodynamics):
In a canonical ensemble, there is no exchange of particles, so the term is zero. Similarly, This gives
where and denote the energies of the reservoir and the system at , respectively. For the second equality we have used the conservation of energy. Substituting into the first equation relating :
which implies, for any state s of the system
where Z is an appropriately chosen "constant" to make total probability 1. (Z is constant provided that the temperature T is invariant.)
where the index s runs through all microstates of the system. Z is sometimes called the Boltzmann sum over states (or "Zustandssumme" in the original German). If we index the summation via the energy eigenvalues instead of all possible states, degeneracy must be taken into account. The probability of our system having energy is simply the sum of the probabilities of all corresponding microstates:
where, with obvious modification,
this is the same result as before.
Comments on this derivation:
Notice that in this formulation, the initial assumption "... suppose the system has total N particles..." is dispensed with. Indeed, the number of particles possessed by the system plays no role in arriving at the distribution. Rather, how many particles would occupy states with energy follows as an easy consequence.
What has been presented above is essentially a derivation of the canonical partition function. As one can see by comparing the definitions, the Boltzmann sum over states is equal to the canonical partition function.
Exactly the same approach can be used to derive Fermi–Dirac and Bose–Einstein statistics. However, there one would replace the canonical ensemble with the grand canonical ensemble, since there is exchange of particles between the system and the reservoir. Also, the system one considers in those cases is a single particle state, not a particle. (In the above discussion, we could have assumed our system to be a single atom.)
| Physical sciences | Statistical mechanics | Physics |
170350 | https://en.wikipedia.org/wiki/Desert%20climate | Desert climate | The desert climate or arid climate (in the Köppen climate classification BWh and BWk) is a dry climate sub-type in which there is a severe excess of evaporation over precipitation. The typically bald, rocky, or sandy surfaces in desert climates are dry and hold little moisture, quickly evaporating the already little rainfall they receive. Covering 14.2% of Earth's land area, hot deserts are the second-most common type of climate on Earth after the Polar climate.
There are two variations of a desert climate according to the Köppen climate classification: a hot desert climate (BWh), and a cold desert climate (BWk). To delineate "hot desert climates" from "cold desert climates", a mean annual temperature of is used as an isotherm so that a location with a BW type climate with the appropriate temperature above this isotherm is classified as "hot arid subtype" (BWh), and a location with the appropriate temperature below the isotherm is classified as "cold arid subtype" (BWk).
Most desert/arid climates receive between of rainfall annually, although some of the most consistently hot areas of Central Australia, the Sahel and Guajira Peninsula can be, due to extreme potential evapotranspiration, classed as arid with the annual rainfall as high as .
Precipitation
Although no part of Earth is known for certain to be rainless, in the Atacama Desert of northern Chile, the average annual rainfall over 17 years was only . Some locations in the Sahara Desert such as Kufra, Libya, record an even drier of rainfall annually. The official weather station in Death Valley, United States reports annually, but in 40 months between 1931 and 1934 a total of just of rainfall was measured.
To determine whether a location has an arid climate, the precipitation threshold is determined. The precipitation threshold (in millimetres) involves first multiplying the average annual temperature in °C by 20, then adding 280 if 70% or more of the total precipitation is in the high-sun summer half of the year (April through September in the Northern Hemisphere, or October through March in the Southern), or 140 if 30–70% of the total precipitation is received during the applicable period, or 0 if less than 30% of the total precipitation is so received there. If the area's annual precipitation is less than half the threshold (50%), it is classified as a BW (desert climate), while 50–100% of the threshold results in a semi-arid climate.
Hot desert climates
Hot desert climates (BWh) are typically found under the subtropical ridge in the lower middle latitudes or the subtropics, often between 20° and 33° north and south latitudes. In these locations, stable descending air and high pressure aloft clear clouds and create hot, arid conditions with intense sunshine. Hot desert climates are found across vast areas of North Africa, West Asia, northwestern parts of the Indian Subcontinent, southwestern Africa, interior Australia, the Southwestern United States, northern Mexico, sections of southeastern Spain, the coast of Peru, and Chile. This makes hot deserts present in every continent except Antarctica.
At the time of high sun (summer), scorching, desiccating heat prevails. Hot-month average temperatures are normally between , and midday readings of are common. The world's absolute heat records, over , are generally in the hot deserts, where the heat potential can be the highest on the planet. This includes the record of in Death Valley, which is currently considered the highest temperature recorded on Earth. Some deserts in the tropics consistently experience very high temperatures all year long, even during wintertime. These locations feature some of the highest annual average temperatures recorded on Earth, exceeding , up to nearly in Dallol, Ethiopia. This last feature is seen in sections of Africa and Arabia. During colder periods of the year, night-time temperatures can drop to freezing or below due to the exceptional radiation loss under the clear skies. However, temperatures rarely drop far below freezing under the hot subtype.
Hot desert climates can be found in the deserts of North Africa such as the wide Sahara Desert, the Libyan Desert or the Nubian Desert; deserts of the Horn of Africa such as the Danakil Desert or the Grand Bara Desert; deserts of Southern Africa such as the Namib Desert or the Kalahari Desert; deserts of West Asia such as the Arabian Desert, or the Syrian Desert; deserts of South Asia such as Dasht-e Lut and Dasht-e Kavir of Iran or the Thar Desert of India and Pakistan; deserts of the United States and Mexico such as the Mojave Desert, the Sonoran Desert or the Chihuahuan Desert; deserts of Australia such as the Simpson Desert or the Great Victoria Desert and many other regions. In Europe, the hot desert climate can only be found on southeastern coast of Spain as well as small inland parts of southeastern, especially parts of the Tabernas Desert.
Hot deserts are lands of extremes: most of them are among the hottest, the driest, and the sunniest places on Earth because of nearly constant high pressure; the almost permanent removal of low-pressure systems, dynamic fronts, and atmospheric disturbances; sinking air motion; dry atmosphere near the surface and aloft; the exacerbated exposure to the sun where solar angles are always high makes this desert inhospitable to most species.
Cold desert climates
Cold desert climates (BWk) usually feature hot (or warm in a few instances), dry summers, though summers are not typically as hot as hot desert climates. Unlike hot desert climates, cold desert climates tend to feature cold, dry winters. Snow tends to be rare in regions with this climate. The Gobi Desert in northern China and Mongolia is one example of a cold desert. Though hot in the summer, it shares the freezing winters of the rest of Inner Asia. Summers in South America's Atacama Desert are mild, with only slight temperature variations between seasons. Cold desert climates are typically found at higher altitudes than hot desert climates and are usually drier than hot desert climates.
Cold desert climates are typically located in temperate zones in the 30s and 40s latitudes, usually in the leeward rain shadow of high mountains, restricting precipitation from the westerly winds. An example of this is the Patagonian Desert in Argentina, bounded by the Andes ranges to its west. In the case of Central Asia, mountains restrict precipitation from the eastern monsoon. The Kyzyl Kum, Taklamakan and Katpana Desert deserts of Central Asia are other significant examples of BWk climates. The Ladakh region and the city of Leh in the Great Himalayas in India also have a cold desert climate. In North America, the cold desert climate occurs in the drier parts of the Great Basin Desert and the Bighorn Basin in Big Horn and Washakie County in Wyoming. The Hautes Plaines, located in the northeastern section of Morocco and in Algeria, is another prominent example of a cold desert climate. In Europe, this climate only occurs in some inland parts of southeastern Spain, such as in Lorca.
Polar climate desert areas in the Arctic and Antarctic regions receive very little precipitation during the year owing to the cold, dry air freezing most precipitation. Polar desert climates have desert-like features that occur in cold desert climates, including intermittent streams, hypersaline lakes, and extremely barren terrain in unglaciated areas such as the McMurdo Dry Valleys of Antarctica. These areas are generally classified as having polar climates because they have average summer temperatures below even if they have some characteristics of extreme non-polar deserts.
Climate charts
Hot deserts
Cold deserts
| Physical sciences | Climates | Earth science |
170396 | https://en.wikipedia.org/wiki/Bark%20%28botany%29 | Bark (botany) | Bark is the outermost layer of stems and roots of woody plants. Plants with bark include trees, woody vines, and shrubs. Bark refers to all the tissues outside the vascular cambium and is a nontechnical term. It overlays the wood and consists of the inner bark and the outer bark. The inner bark, which in older stems is living tissue, includes the innermost layer of the periderm. The outer bark on older stems includes the dead tissue on the surface of the stems, along with parts of the outermost periderm and all the tissues on the outer side of the periderm. The outer bark on trees which lies external to the living periderm is also called the rhytidome.
Products derived from bark include bark shingle siding and wall coverings, spices, and other flavorings, tanbark for tannin, resin, latex, medicines, poisons, various hallucinogenic chemicals, and cork. Bark has been used to make cloth, canoes, and ropes and used as a surface for paintings and map making. A number of plants are also grown for their attractive or interesting bark colorations and surface textures or their bark is used as landscape mulch.
The process of removing bark is decortication and a log or trunk from which bark has been removed is said to be decorticated.
Botanical description
Bark is present only on woody plants - herbaceous plants and stems of young plants lack bark.
From the outside to the inside of a mature woody stem, the layers include the following:
Bark
Periderm
Cork (phellem or suber), includes the rhytidome
Cork cambium (phellogen)
Phelloderm
Cortex
Phloem
Vascular cambium
Wood (xylem)
Sapwood (alburnum)
Heartwood (duramen)
Pith (medulla)
In young stems, which lack what is commonly called bark, the tissues are, from the outside to the inside:
Epidermis, which may be replaced by periderm
Cortex
Primary and secondary phloem
Vascular cambium
Secondary and primary xylem.
Cork cell walls contain suberin, a waxy substance which protects the stem against water loss, the invasion of insects into the stem, and prevents infections by bacteria and fungal spores. The cambium tissues, i.e., the cork cambium and the vascular cambium, are the only parts of a woody stem where cell division occurs; undifferentiated cells in the vascular cambium divide rapidly to produce secondary xylem to the inside and secondary phloem to the outside. Phloem is a nutrient-conducting tissue composed of sieve tubes or sieve cells mixed with parenchyma and fibers. The cortex is the primary tissue of stems and roots. In stems the cortex is between the epidermis layer and the phloem, in roots the inner layer is not phloem but the pericycle.
As the stem ages and grows, changes occur that transform the surface of the stem into the bark. The epidermis is a layer of cells that cover the plant body, including the stems, leaves, flowers and fruits, that protects the plant from the outside world. In old stems the epidermal layer, cortex, and primary phloem become separated from the inner tissues by thicker formations of cork. Due to the thickening cork layer these cells die because they do not receive water and nutrients. This dead layer is the rough corky bark that forms around tree trunks and other stems.
Cork, sometimes confused with bark in colloquial speech, is the outermost layer of a woody stem, derived from the cork cambium. It serves as protection against damage from parasites, herbivorous animals and diseases, as well as dehydration and fire.
Periderm
Often a secondary covering called the periderm forms on small woody stems and many non-woody plants, which is composed of cork (phellem), the cork cambium (phellogen), and the phelloderm. The periderm forms from the phellogen which serves as a lateral meristem. The periderm replaces the epidermis, and acts as a protective covering like the epidermis. Mature phellem cells have suberin in their walls to protect the stem from desiccation and pathogen attack. Older phellem cells are dead, as is the case with woody stems. The skin on the potato tuber (which is an underground stem) constitutes the cork of the periderm.
In woody plants, the epidermis of newly grown stems is replaced by the periderm later in the year. As the stems grow a layer of cells form under the epidermis, called the cork cambium, these cells produce cork cells that turn into cork. A limited number of cell layers may form interior to the cork cambium, called the phelloderm.
As the stem grows, the cork cambium produces new layers of cork which are impermeable to gases and water and the cells outside the periderm, namely the epidermis, cortex and older secondary phloem die.
Within the periderm are lenticels, which form during the production of the first periderm layer. Since there are living cells within the cambium layers that need to exchange gases during metabolism, these lenticels, because they have numerous intercellular spaces, allow gaseous exchange with the outside atmosphere. As the bark develops, new lenticels are formed within the cracks of the cork layers.
Rhytidome
The rhytidome is the most familiar part of bark, being the outer layer that covers the trunks of trees. It is composed mostly of dead cells and is produced by the formation of multiple layers of suberized periderm, cortical and phloem tissue. The rhytidome is especially well developed in older stems and roots of trees. In shrubs, older bark is quickly exfoliated and thick rhytidome accumulates. It is generally thickest and most distinctive at the trunk or bole (the area from the ground to where the main branching starts) of the tree.
Chemical composition
Bark tissues make up by weight between 10 and 20% of woody vascular plants and consists of various biopolymers, tannins, lignin, suberin and polysaccharides. Up to 40% of the bark tissue is made of lignin, which forms an important part of a plant, providing structural support by crosslinking between different polysaccharides, such as cellulose.
Condensed tannin, which is in fairly high concentration in bark tissue, is thought to inhibit decomposition. It could be due to this factor that the degradation of lignin is far less pronounced in bark tissue than it is in wood.
It has been proposed that, in the cork layer (the phellogen), suberin acts as a barrier to microbial degradation and so protects the internal structure of the plant.
Analysis of the lignin in the bark wall during decay by the white-rot fungi Lentinula edodes (Shiitake mushroom) using 13C NMR revealed that the lignin polymers contained more Guaiacyl lignin units than Syringyl units compared to the interior of the plant. Guaiacyl units are less susceptible to degradation as, compared to syringyl, they contain fewer aryl-aryl bonds, can form a condensed lignin structure, and have a lower redox potential. This could mean that the concentration and type of lignin units could provide additional resistance to fungal decay for plants protected by bark.
Damage and repair
Bark can sustain damage from environmental factors, such as frost crack and sun scald, as well as biological factors, such as woodpecker and boring beetle attacks. Male deer and other male members of the Cervidae (deer family) can cause extensive bark damage during the rutting season by rubbing their antlers against the tree to remove their velvet.
The bark is often damaged by being bound to stakes or wrapped with wires. In the past, this damage was called bark-galling and was treated by applying clay laid on the galled place and binding it up with hay. In modern usage, "galling" most typically refers to a type of abnormal growth on a plant caused by insects or pathogens.
Bark damage can have several detrimental effects on the plant. Bark serves as a physical barrier to disease pressure, especially from fungi, so its removal makes the plant more susceptible to disease. Damage or destruction of the phloem impedes the transport of photosynthetic products throughout the plant; in extreme cases, when a band of phloem all the way around the stem is removed, the plant will usually quickly die. Bark damage in horticultural applications, as in gardening and public landscaping, results in often unwanted aesthetic damage.
The degree to which woody plants are able to repair gross physical damage to their bark is quite variable across species and type of damage. Some are able to produce a callus growth which heals over the wound rapidly, but leaves a clear scar, whilst others such as oaks do not produce an extensive callus repair. Sap is sometimes produced to seal the damaged area against disease and insect intrusion.
A number of living organisms live in or on bark, including insects, fungi and other plants like mosses, algae and other vascular plants. Many of these organisms are pathogens or parasites but some also have symbiotic relationships.
Uses
The inner bark (phloem) of some trees is edible. In hunter-gatherer societies and in times of famine, it is harvested and used as a food source. In Scandinavia, bark bread is made from rye to which the toasted and ground innermost layer of bark of scots pine or birch is added. The Sami people of far northern Europe use large sheets of Pinus sylvestris bark that are removed in the spring, prepared and stored for use as a staple food resource. The inner bark is eaten fresh, dried or roasted.
Bark can be used as a construction material, and was used widely in pre-industrial societies. Some barks, particularly Birch bark, can be removed in long sheets and other mechanically cohesive structures, allowing the bark to be used in the construction of canoes, as the drainage layer in roofs, for shoes, backpacks, and other useful items. Bark was also used as a construction material in settler colonial societies, particularly Australia, both as exterior wall cladding and as a roofing material.
In the cork oak (Quercus suber) the bark is thick enough to be harvested as a cork product without killing the tree; in this species the bark may get very thick (e.g. more than 20 cm has been reported).
Some s have significantly different phytochemical content from other parts. Some of these phytochemicals have pesticidal, culinary, or medicinally and culturally important ethnopharmacological properties.
Bark contains strong fibres known as bast, and there is a long tradition in northern Europe of using bark from coppiced young branches of the small-leaved lime (Tilia cordata) to produce cordage and rope, used for example in the rigging of Viking Age longships.
Among the commercial products made from bark are cork, cinnamon, quinine (from the bark of Cinchona) and aspirin (from the bark of willow trees). The bark of some trees, notably oak (Quercus robur) is a source of tannic acid, which is used in tanning. Bark chips generated as a by-product of lumber production are often used in bark mulch. Bark is important to the horticultural industry since in shredded form it is used for plants that do not thrive in ordinary soil, such as epiphytes.
Wood bark contains lignin which when pyrolyzed yields a liquid bio-oil product rich in natural phenol derivatives. These are used as a replacement for fossil-based phenols in phenol-formaldehyde (PF) resins used in Oriented Strand Board (OSB) and plywood.
Gallery
| Biology and health sciences | Plant stem | null |
170406 | https://en.wikipedia.org/wiki/Silt | Silt | Silt is granular material of a size between sand and clay and composed mostly of broken grains of quartz. Silt may occur as a soil (often mixed with sand or clay) or as sediment mixed in suspension with water. Silt usually has a floury feel when dry, and lacks plasticity when wet. Silt can also be felt by the tongue as granular when placed on the front teeth (even when mixed with clay particles).
Silt is a common material, making up 45% of average modern mud. It is found in many river deltas and as wind-deposited accumulations, particularly in central Asia, north China, and North America. It is produced in both very hot climates (through such processes as collisions of quartz grains in dust storms) and very cold climates (through such processes as glacial grinding of quartz grains.)
Loess is soil rich in silt which makes up some of the most fertile agricultural land on Earth. However, silt is very vulnerable to erosion, and it has poor mechanical properties, making construction on silty soil problematic. The failure of the Teton Dam in 1976 has been attributed to the use of unsuitable loess in the dam core, and liquefication of silty soil is a significant earthquake hazard. Windblown and waterborne silt are significant forms of environmental pollution, often exacerbated by poor farming practices.
Description
Silt is detritus (fragments of weathered and eroded rock) with properties intermediate between sand and clay. A more precise definition of silt used by geologists is that it is detrital particles with sizes between 1/256 and 1/16 mm (about 4 to 63 microns). This corresponds to particles between 8 and 4 phi units on the Krumbein phi scale. Other geologists define silt as detrital particles between 2 and 63 microns or 9 to 4 phi units. A third definition is that silt is fine-grained detrital material composed of quartz rather than clay minerals. Since most clay mineral particles are smaller than 2 microns, while most detrital particles between 2 and 63 microns in size are composed of broken quartz grains, there is good agreement between these definitions in practice.
The upper size limit of 1/16 mm or 63 microns corresponds to the smallest particles that can be discerned with the unaided eye. It also corresponds to a Tanner gap in the distribution of particle sizes in sediments: Particles between 120 and 30 microns in size are scarce in most sediments, suggesting that the distinction between sand and silt has physical significance. As noted above, the lower limit of 2 to 4 microns corresponds to the transition from particles that are predominantly broken quartz grains to particles that are predominantly clay mineral particles.
Assallay and coinvestigators further divide silt into three size ranges: C (2–5 microns), which represents post-glacial clays and desert dust; D1 (20–30 microns) representing "traditional" loess; and D2 (60 microns) representing the very coarse North African loess.
Silt can be distinguished from clay in the field by its lack of plasticity or cohesiveness and by its grain size. Silt grains are large enough to give silt a gritty feel, particularly if a sample is placed between the teeth. Clay-size particles feel smooth between the teeth. The proportions of coarse and fine silt in a sediment sample are determined more precisely in the laboratory using the pipette method, which is based on settling rate via Stokes' law and gives the particle size distribution accordingly. The mineral composition of silt particles can be determined with a petrographic microscope for grain sizes as low as 10 microns.
Vadose silt is silt-sized calcite crystals found in pore spaces and vugs in limestone. This is emplaced as sediment is carried through the vadose zone to be deposited in pore space.
Definitions
ASTM American Standard of Testing Materials: 200 sieve – 0.005 mm.
USDA United States Department of Agriculture 0.05–0.002 mm.
ISSS International Society of Soil Science 0.02–0.002 mm.
Civil engineers in the United States define silt as material made of particles that pass a number 200 sieve (0.074 mm or less) but show little plasticity when wet and little cohesion when air-dried. The International Society of Soil Science (ISSS) defines silt as soil containing 80% or more of particles between 0.002 mm to 0.02 mm in size while the U.S. Department of Agriculture puts the cutoff at 0.05mm. The term silt is also used informally for material containing much sand and clay as well as silt-sized particles, or for mud suspended in water.
Occurrence
Silt is a very common material, and it has been estimated that there are a billion trillion trillion (1033) silt grains worldwide. Silt is abundant in eolian and alluvial deposits, including river deltas, such as the Nile and Niger River deltas. Bangladesh is largely underlain by silt deposits of the Ganges delta. Silt is also abundant in northern China, central Asia, and North America. However, silt is relatively uncommon in the tropical regions of the world.
Silt is commonly found in suspension in river water, and it makes up over 0.2% of river sand. It is abundant in the matrix between the larger sand grains of graywackes. Modern mud has an average silt content of 45%. Silt is often found in mudrock as thin laminae, as clumps, or dispersed throughout the rock. Laminae suggest deposition in a weak current that winnows the silt of clay, while clumps suggest an origin as fecal pellets. Where silt is dispersed throughout the mudrock, it likely was deposited by rapid processes, such as flocculation. Sedimentary rock composed mainly of silt is known as siltstone.
Silt is common throughout the geologic record, but it seems to be particularly common in Quaternary formations. This may be because deposition of silt is favored by the glaciation and arctic conditions characteristic of the Quaternary. Silt is sometimes known as rock flour or glacier meal, especially when produced by glacial action. Silt suspended in water draining from glaciers is sometimes known as rock milk or moonmilk.
Sources
A simple explanation for silt formation is that it is a straightforward continuation to a smaller scale of the disintegration of rock into gravel and sand. However, the presence of a Tanner gap between sand and silt (a scarcity of particles with sizes between 30 and 120 microns) suggests that different physical processes produce sand and silt. The mechanisms of silt formation have been studied extensively in the laboratory and compared with field observations. These show that silt formation requires high-energy processes acting over long periods of time, but such processes are present in diverse geologic settings.
Quartz silt grains are usually found to have a platy or bladed shape. This may be characteristic of how larger grains abrade, or reflect the shape of small quartz grains in foliated metamorphic rock, or arise from authigenic growth of quartz grains parallel to bedding in sedimentary rock. Theoretically, particles formed by random fracturing of an isotropic material, such as quartz, naturally tend to be blade-shaped. The size of silt grains produced by abrasion or shattering of larger grains may reflect defects in the crystal structure of the quartz, known as Moss defects. Such defects are produced by tectonic deformation of the parent rock, and also arise from the high-low transition of quartz: Quartz experiences a sharp decrease in volume when it cools below a temperature of about , which creates strain and crystal defects in the quartz grains in a cooling body of granite.
Mechanisms for silt production include:
Erosion of initially silt-sized grains from low-grade metamorphic rock.
Production of silt-sized grains from fracture of larger grains during initial rock weathering and soil formation, through processes such as frost shattering and haloclasty. This produces silt particles whose size of 10–30 microns is determined by Moss defects.
Production of silt-sized grains from grain-to-grain impact during transport of coarser sediments.
Formation of authigenic quartz during weathering to clay.
Crystallization of the tests of siliceous organisms deposited in mudrock.
Laboratory experiments have produced contradictory results regarding the effectiveness of various silt production mechanisms. This may be due to the use of vein or pegmatite quartz in some of the experiments. Both materials form under conditions promoting ideal crystal growth, and may lack the Moss defects of quartz grains in granites. Thus production of silt from vein quartz is very difficult by any mechanism, whereas production of silt from granite quartz proceeds readily by any of a number of mechanisms. However, the main process is likely abrasion through transport, including fluvial comminution, aeolian attrition and glacial grinding.
Because silt deposits (such as loess, a soil composed mostly of silt) seem to be associated with glaciated or mountainous regions in Asia and North America, much emphasis has been placed on glacial grinding as a source of silt. High Asia has been identified as a major generator of silt, which accumulated to form the fertile soils of north India and Bangladesh, and the loess of central Asia and north China. Loess has long been thought to be absent or rare in deserts lacking nearby mountains (Sahara, Australia). However, laboratory experiments show eolian and fluvial processes can be quite efficient at producing silt, as can weathering in tropical climates. Silt seems to be produced in great quantities in dust storms, and silt deposits found in Israel, Tunisia, Nigeria, and Saudi Arabia cannot be attributed to glaciation. Furthermore, desert source areas in Asia may be more important for loess formation than previously thought. Part of the problem may be the conflation of high rates of production with environments conducive to deposition and preservation, which favors glacial climates more than deserts.
Loess associated with glaciation and cold weathering may be distinguishable from loess associated with hot regions by the size distribution. Glacial loess has a typical particle size of about 25 microns. Desert loess contains either larger or smaller particles, with the fine silt produced in dust storms and the coarse silt fraction possibly representing the fine particle tail of sand production.
Human impact
Loess underlies some of the most productive agricultural land worldwide. However, it is very susceptible to erosion. The quartz particles in silt do not themselves provide nutrients, but they promote excellent soil structure, and silt-sized particles of other minerals, present in smaller amounts, provide the necessary nutrients. Silt, deposited by annual floods along the Nile River, created the rich, fertile soil that sustained the Ancient Egyptian civilization. The closure of the Aswan High Dam has cut off this source of silt, and the fertility of the Nile delta is deteriorating.
Loess tends to lose strength when wetted, and this can lead to failure of building foundations. The silty material has an open structure that collapses when wet. Quick clay (a combination of very fine silt and clay-sized particles from glacial grinding) is a particular challenge for civil engineering.
The failure of the Teton Dam has been attributed to the use of loess from the Snake River floodplain in the core of the dam. Loess lacks the necessary plasticity for use in a dam core, but its properties were poorly understood, even by the U.S. Bureau of Reclamation, with its wealth of experience building earthen dams.
Silt is susceptible to liquefaction during strong earthquakes due to its lack of plasticity. This has raised concerns about the earthquake damage potential in the silty soil of the central United States in the event of a major earthquake in the New Madrid Seismic Zone.
Environmental impacts
Silt is easily transported in water and is fine enough to be carried long distances by air in the form of dust. While the coarsest silt particles (60 micron) settle out of a meter of still water in just five minutes, the finest silt grains (2 microns) can take several days to settle out of still water. When silt appears as a pollutant in water the phenomenon is known as siltation.
Silt deposited by the Mississippi River throughout the 20th century has decreased due to a system of levees, contributing to the disappearance of protective wetlands and barrier islands in the delta region surrounding New Orleans.
In southeast Bangladesh, in the Noakhali district, cross dams were built in the 1960s whereby silt gradually started forming new land called "chars". The district of Noakhali has gained more than of land in the past 50 years. With Dutch funding, the Bangladeshi government began to help develop older chars in the late 1970s, and the effort has since become a multi-agency operation building roads, culverts, embankments, cyclone shelters, toilets and ponds, as well as distributing land to settlers. By fall 2010, the program will have allotted some to 21,000 families.
A main source of silt in urban rivers is disturbance of soil by construction activity. A main source in rural rivers is erosion from plowing of farm fields, clearcutting or slash and burn treatment of forests.
Culture
The fertile black silt of the Nile river's banks is a symbol of rebirth, associated with the Egyptian god Anubis.
| Physical sciences | Sedimentology | Earth science |
170417 | https://en.wikipedia.org/wiki/T%20cell | T cell | T cells are one of the important types of white blood cells of the immune system and play a central role in the adaptive immune response. T cells can be distinguished from other lymphocytes by the presence of a T-cell receptor (TCR) on their cell surface.
T cells are born from hematopoietic stem cells, found in the bone marrow. Developing T cells then migrate to the thymus gland to develop (or mature). T cells derive their name from the thymus. After migration to the thymus, the precursor cells mature into several distinct types of T cells. T cell differentiation also continues after they have left the thymus. Groups of specific, differentiated T cell subtypes have a variety of important functions in controlling and shaping the immune response.
One of these functions is immune-mediated cell death, and it is carried out by two major subtypes: CD8+ "killer" (cytotoxic) and CD4+ "helper" T cells. (These are named for the presence of the cell surface proteins CD8 or CD4.) CD8+ T cells, also known as "killer T cells", are cytotoxic – this means that they are able to directly kill virus-infected cells, as well as cancer cells. CD8+ T cells are also able to use small signalling proteins, known as cytokines, to recruit other types of cells when mounting an immune response. A different population of T cells, the CD4+ T cells, function as "helper cells". Unlike CD8+ killer T cells, the CD4+ helper T (TH) cells function by further activating memory B cells and cytotoxic T cells, which leads to a larger immune response. The specific adaptive immune response regulated by the TH cell depends on its subtype (such as T-helper1, T-helper2, T-helper17, regulatory T-cell), which is distinguished by the types of cytokines they secrete.
Regulatory T cells are yet another distinct population of T cells that provide the critical mechanism of tolerance, whereby immune cells are able to distinguish invading cells from "self". This prevents immune cells from inappropriately reacting against one's own cells, known as an "autoimmune" response. For this reason, these regulatory T cells have also been called "suppressor" T cells. These same regulatory T cells can also be co-opted by cancer cells to prevent the recognition of, and an immune response against, tumor cells.
Development
Origin, early development and migration to the thymus
All T cells originate from c-kit+Sca1+ haematopoietic stem cells (HSC) which reside in the bone marrow. In some cases, the origin might be the foetal liver during embryonic development. The HSC then differentiate into multipotent progenitors (MPP) which retain the potential to become both myeloid and lymphoid cells. The process of differentiation then proceeds to a common lymphoid progenitor (CLP), which can only differentiate into T, B or NK cells. These CLP cells then migrate via the blood to the thymus, where they engraft:. Henceforth they are known as thymocytes, the immature stage of a T cell.
The earliest cells which arrived in the thymus are commonly termed double-negative, as they express neither the CD4 nor CD8 co-receptor. The newly arrived CLP cells are CD4−CD8−CD44+CD25−ckit+ cells, and are termed early thymic progenitor (ETP) cells. These cells will then undergo a round of division and downregulate c-kit and are termed double-negative one (DN1) cells. To become T cells, the thymocytes must undergo multiple DN stages as well as positive selection and negative selection.
Double negative thymocytes can be identified by the surface expression of CD2, CD5 and CD7. Still during the double negative stages, CD34 expression stops and CD1 is expressed. Expression of both CD4 and CD8 makes them double positive, and matures into either CD4+ or CD8+ cells.
TCR development
A critical step in T cell maturation is making a functional T cell receptor (TCR). Each mature T cell will ultimately contain a unique TCR that reacts to a random pattern, allowing the immune system to recognize many different types of pathogens. This process is essential in developing immunity to threats that the immune system has not encountered before, since due to random variation there will always be at least one TCR to match any new pathogen.
A thymocyte can only become an active T cell when it survives the process of developing a functional TCR. The TCR consists of two major components, the alpha and beta chains. These both contain random elements designed to produce a wide variety of different TCRs, but due to this huge variety they must be tested to make sure they work at all. First, the thymocytes attempt to create a functional beta chain, testing it against a 'mock' alpha chain. Then they attempt to create a functional alpha chain. Once a working TCR has been produced, the cells then must test if their TCR will identify threats correctly, and to do this it is required to recognize the body’s major histocompatibility complex (MHC) in a process known as positive selection. The thymocyte must also ensure that it does not react adversely to "self" antigens, called negative selection. If both positive and negative selection are successful, the TCR becomes fully operational and the thymocyte becomes a T cell.
TCR β-chain selection
At the DN2 stage (CD44+CD25+), cells upregulate the recombination genes RAG1 and RAG2 and re-arrange the TCRβ locus, combining V-D-J recombination and constant region genes in an attempt to create a functional TCRβ chain. As the developing thymocyte progresses through to the DN3 stage (CD44−CD25+), the thymocyte expresses an invariant α-chain called pre-Tα alongside the TCRβ gene. If the rearranged β-chain successfully pairs with the invariant α-chain, signals are produced which cease rearrangement of the β-chain (and silence the alternate allele). Although these signals require the pre-TCR at the cell surface, they are independent of ligand binding to the pre-TCR. If the chains successfully pair a pre-TCR forms, and the cell downregulates CD25 and is termed a DN4 cell (CD25−CD44−). These cells then undergo a round of proliferation, and begin to re-arrange the TCRα locus during the double-positive stage.
Positive selection
The process of positive selection takes 3 to 4 days and occurs in the thymic cortex. Double-positive thymocytes (CD4+/CD8+) migrate deep into the thymic cortex, where they are presented with self-antigens. These self-antigens are expressed by thymic cortical epithelial cells on MHC molecules, which reside on the surface of cortical epithelial cells. Only thymocytes that interact well with MHC-I or MHC-II will receive a vital "survival signal", while those that cannot interact strongly enough will receive no signal and die from neglect. This process ensures that the surviving thymocytes will have an 'MHC affinity' that means they will exhibit stronger binding affinity for specific MHC alleles in that organism. The vast majority of developing thymocytes will not pass positive selection, and die during this process.
A thymocyte's fate is determined during positive selection. Double-positive cells (CD4+/CD8+) that interact well with MHC class II molecules will eventually become CD4+ "helper" cells, whereas thymocytes that interact well with MHC class I molecules mature into CD8+ "killer" cells. A thymocyte becomes a CD4+ cell by down-regulating expression of its CD8 cell surface receptors. If the cell does not lose its signal, it will continue downregulating CD8 and become a CD4+, both CD8+ and CD4+ cells are now single positive cells.
This process does not filter for thymocytes that may cause autoimmunity. The potentially autoimmune cells are removed by the following process of negative selection, which occurs in the thymic medulla.
Negative selection
Negative selection removes thymocytes that are capable of strongly binding with "self" MHC molecules. Thymocytes that survive positive selection migrate towards the boundary of the cortex and medulla in the thymus. While in the medulla, they are again presented with a self-antigen presented on the MHC complex of medullary thymic epithelial cells (mTECs). mTECs must be Autoimmune regulator positive (AIRE+) to properly express tissue-specific antigens on their MHC class I peptides. Some mTECs are phagocytosed by thymic dendritic cells; this makes them AIRE− antigen presenting cells (APCs), allowing for presentation of self-antigens on MHC class II molecules (positively selected CD4+ cells must interact with these MHC class II molecules, thus APCs, which possess MHC class II, must be present for CD4+ T-cell negative selection). Thymocytes that interact too strongly with the self-antigen receive an apoptotic signal that leads to cell death. However, some of these cells are selected to become Treg cells. The remaining cells exit the thymus as mature naive T cells, also known as recent thymic emigrants. This process is an important component of central tolerance and serves to prevent the formation of self-reactive T cells that are capable of inducing autoimmune diseases in the host.
TCR development summary
β-selection is the first checkpoint, where thymocytes that are able to form a functional pre-TCR (with an invariant alpha chain and a functional beta chain) are allowed to continue development in the thymus. Next, positive selection checks that thymocytes have successfully rearranged their TCRα locus and are capable of recognizing MHC molecules with appropriate affinity. Negative selection in the medulla then eliminates thymocytes that bind too strongly to self-antigens expressed on MHC molecules. These selection processes allow for tolerance of self by the immune system. Typical naive T cells that leave the thymus (via the corticomedullary junction) are self-restricted, self-tolerant, and single positive.
Thymic output
About 98% of thymocytes die during the development processes in the thymus by failing either positive selection or negative selection, whereas the other 2% survive and leave the thymus to become mature immunocompetent T cells.
The thymus contributes fewer cells as a person ages. As the thymus shrinks by about 3% a year throughout middle age, a corresponding fall in the thymic production of naive T cells occurs, leaving peripheral T cell expansion and regeneration to play a greater role in protecting older people.
Types of T cell
T cells are grouped into a series of subsets based on their function. CD4 and CD8 T cells are selected in the thymus, but undergo further differentiation in the periphery to specialized cells which have different functions. T cell subsets were initially defined by function, but also have associated gene or protein expression patterns.
Conventional adaptive T cells
Helper CD4+ T cells
T helper cells (TH cells) assist other lymphocytes, including the maturation of B cells into plasma cells and memory B cells, and activation of cytotoxic T cells and macrophages. These cells are also known as CD4+ T cells as they express the CD4 glycoprotein on their surfaces. Helper T cells become activated when they are presented with peptide antigens by MHC class II molecules, which are expressed on the surface of antigen-presenting cells (APCs). Once activated, they divide rapidly and secrete cytokines that regulate or assist the immune response. These cells can differentiate into one of several subtypes, which have different roles. Cytokines direct T cells into particular subtypes.
Cytotoxic CD8+ T cells
Cytotoxic T cells (TC cells, CTLs, T-killer cells, killer T cells) destroy virus-infected cells and tumor cells, and are also implicated in transplant rejection. These cells are defined by the expression of the CD8 protein on their cell surface. Cytotoxic T cells recognize their targets by binding to short peptides (8-11 amino acids in length) associated with MHC class I molecules, present on the surface of all nucleated cells. Cytotoxic T cells also produce the key cytokines IL-2 and IFNγ. These cytokines influence the effector functions of other cells, in particular macrophages and NK cells.
Memory T cells
Antigen-naive T cells expand and differentiate into memory and effector T cells after they encounter their cognate antigen within the context of an MHC molecule on the surface of a professional antigen presenting cell (e.g. a dendritic cell). Appropriate co-stimulation must be present at the time of antigen encounter for this process to occur. Historically, memory T cells were thought to belong to either the effector or central memory subtypes, each with their own distinguishing set of cell surface markers (see below). Subsequently, numerous new populations of memory T cells were discovered including tissue-resident memory T (Trm) cells, stem memory TSCM cells, and virtual memory T cells. The single unifying theme for all memory T cell subtypes is that they are long-lived and can quickly expand to large numbers of effector T cells upon re-exposure to their cognate antigen. By this mechanism they provide the immune system with "memory" against previously encountered pathogens. Memory T cells may be either CD4+ or CD8+ and usually express CD45RO.
Memory T cell subtypes:
Central memory T cells (TCM cells) express CD45RO, C-C chemokine receptor type 7 (CCR7), and L-selectin (CD62L). Central memory T cells also have intermediate to high expression of CD44. This memory subpopulation is commonly found in the lymph nodes and in the peripheral circulation. (Note- CD44 expression is usually used to distinguish murine naive from memory T cells).
Effector memory T cells (TEM cells and TEMRA cells) express CD45RO but lack expression of CCR7 and L-selectin. They also have intermediate to high expression of CD44. These memory T cells lack lymph node-homing receptors and are thus found in the peripheral circulation and tissues. TEMRA stands for terminally differentiated effector memory cells re-expressing CD45RA, which is a marker usually found on naive T cells.
Tissue-resident memory T cells (TRM) occupy tissues (skin, lung, etc.) without recirculating. One cell surface marker that has been associated with TRM is the intern αeβ7, also known as CD103.
Virtual memory T cells (TVM) differ from the other memory subsets in that they do not originate following a strong clonal expansion event. Thus, although this population as a whole is abundant within the peripheral circulation, individual virtual memory T cell clones reside at relatively low frequencies. One theory is that homeostatic proliferation gives rise to this T cell population. Although CD8 virtual memory T cells were the first to be described, it is now known that CD4 virtual memory cells also exist.
Regulatory CD4+ T cells
Regulatory T cells are crucial for the maintenance of immunological tolerance. Their major role is to shut down T cell–mediated immunity toward the end of an immune reaction and to suppress autoreactive T cells that escaped the process of negative selection in the thymus.
Two major classes of CD4+ Treg cells have been described—FOXP3+ Treg cells and FOXP3− Treg cells.
Regulatory T cells can develop either during normal development in the thymus, and are then known as thymic Treg cells, or can be induced peripherally and are called peripherally derived Treg cells. These two subsets were previously called "naturally occurring" and "adaptive" (or "induced"), respectively. Both subsets require the expression of the transcription factor FOXP3 which can be used to identify the cells. Mutations of the FOXP3 gene can prevent regulatory T cell development, causing the fatal autoimmune disease IPEX.
Several other types of T cells have suppressive activity, but do not express FOXP3 constitutively. These include Tr1 and Th3 cells, which are thought to originate during an immune response and act by producing suppressive molecules. Tr1 cells are associated with IL-10, and Th3 cells are associated with TGF-beta. Recently, Th17 cells have been added to this list.
Innate-like T cells
Innate-like T cells or unconventional T cells represent some subsets of T cells that behave differently in immunity. They trigger rapid immune responses, regardless of the major histocompatibility complex (MHC) expression, unlike their conventional counterparts (CD4 T helper cells and CD8 cytotoxic T cells), which are dependent on the recognition of peptide antigens in the context of the MHC molecule. Overall, there are three large populations of unconventional T cells: NKT cells, MAIT cells, and gammadelta T cells. Now, their functional roles are already being well established in the context of infections and cancer. Furthermore, these T cell subsets are being translated into many therapies against malignancies such as leukemia, for example.
Natural killer T cell
Natural killer T cells (NKT cells – not to be confused with natural killer cells of the innate immune system) bridge the adaptive immune system with the innate immune system. Unlike conventional T cells that recognize protein peptide antigens presented by major histocompatibility complex (MHC) molecules, NKT cells recognize glycolipid antigens presented by CD1d. Once activated, these cells can perform functions ascribed to both helper and cytotoxic T cells: cytokine production and release of cytolytic/cell killing molecules. They are also able to recognize and eliminate some tumor cells and cells infected with herpes viruses.
Mucosal associated invariant T cells
Mucosal associated invariant T (MAIT) cells display innate, effector-like qualities. In humans, MAIT cells are found in the blood, liver, lungs, and mucosa, defending against microbial activity and infection. The MHC class I-like protein, MR1, is responsible for presenting bacterially-produced vitamin B metabolites to MAIT cells. After the presentation of foreign antigen by MR1, MAIT cells secrete pro-inflammatory cytokines and are capable of lysing bacterially-infected cells. MAIT cells can also be activated through MR1-independent signaling. In addition to possessing innate-like functions, this T cell subset supports the adaptive immune response and has a memory-like phenotype. Furthermore, MAIT cells are thought to play a role in autoimmune diseases, such as multiple sclerosis, arthritis and inflammatory bowel disease, although definitive evidence is yet to be published.
Gamma delta T cells
Gamma delta T cells (γδ T cells) represent a small subset of T cells which possess a γδ TCR rather than the αβ TCR on the cell surface. The majority of T cells express αβ TCR chains. This group of T cells is much less common in humans and mice (about 2% of total T cells) and are found mostly in the gut mucosa, within a population of intraepithelial lymphocytes. In rabbits, sheep, and chickens, the number of γδ T cells can be as high as 60% of total T cells. The antigenic molecules that activate γδ T cells are still mostly unknown. However, γδ T cells are not MHC-restricted and seem to be able to recognize whole proteins rather than requiring peptides to be presented by MHC molecules on APCs. Some murine γδ T cells recognize MHC class IB molecules. Human γδ T cells that use the Vγ9 and Vδ2 gene fragments constitute the major γδ T cell population in peripheral blood. These cells are unique in that they specifically and rapidly respond to a set of nonpeptidic phosphorylated isoprenoid precursors, collectively named phosphoantigens, which are produced by virtually all living cells. The most common phosphoantigens from animal and human cells (including cancer cells) are isopentenyl pyrophosphate (IPP) and its isomer dimethylallyl pyrophosphate (DMPP). Many microbes produce the active compound hydroxy-DMAPP (HMB-PP) and corresponding mononucleotide conjugates, in addition to IPP and DMAPP. Plant cells produce both types of phosphoantigens. Drugs activating human Vγ9/Vδ2 T cells comprise synthetic phosphoantigens and aminobisphosphonates, which upregulate endogenous IPP/DMAPP.
Activation
Activation of CD4+ T cells occurs through the simultaneous engagement of the T-cell receptor and a co-stimulatory molecule (like CD28, or ICOS) on the T cell by the major histocompatibility complex (MHCII) peptide and co-stimulatory molecules on the APC. Both are required for production of an effective immune response; in the absence of co-stimulation, T cell receptor signalling alone results in anergy. The signalling pathways downstream from co-stimulatory molecules usually engages the PI3K pathway generating PIP3 at the plasma membrane and recruiting PH domain containing signaling molecules like PDK1 that are essential for the activation of PKC-θ, and eventual IL-2 production. Optimal CD8+ T cell response relies on CD4+ signalling. CD4+ cells are useful in the initial antigenic activation of naive CD8 T cells, and sustaining memory CD8+ T cells in the aftermath of an acute infection. Therefore, activation of CD4+ T cells can be beneficial to the action of CD8+ T cells.
The first signal is provided by binding of the T cell receptor to its cognate peptide presented on MHCII on an APC. MHCII is restricted to so-called professional antigen-presenting cells, like dendritic cells, B cells, and macrophages, to name a few. The peptides presented to CD8+ T cells by MHC class I molecules are 8–13 amino acids in length; the peptides presented to CD4+ cells by MHC class II molecules are longer, usually 12–25 amino acids in length, as the ends of the binding cleft of the MHC class II molecule are open.
The second signal comes from co-stimulation, in which surface receptors on the APC are induced by a relatively small number of stimuli, usually products of pathogens, but sometimes breakdown products of cells, such as necrotic-bodies or heat shock proteins. The only co-stimulatory receptor expressed constitutively by naive T cells is CD28, so co-stimulation for these cells comes from the CD80 and CD86 proteins, which together constitute the B7 protein, (B7.1 and B7.2, respectively) on the APC. Other receptors are expressed upon activation of the T cell, such as OX40 and ICOS, but these largely depend upon CD28 for their expression. The second signal licenses the T cell to respond to an antigen. Without it, the T cell becomes anergic, and it becomes more difficult for it to activate in future. This mechanism prevents inappropriate responses to self, as self-peptides will not usually be presented with suitable co-stimulation. Once a T cell has been appropriately activated (i.e. has received signal one and signal two) it alters its cell surface expression of a variety of proteins. Markers of T cell activation include CD69, CD71 and CD25 (also a marker for Treg cells), and HLA-DR (a marker of human T cell activation). CTLA-4 expression is also up-regulated on activated T cells, which in turn outcompetes CD28 for binding to the B7 proteins. This is a checkpoint mechanism to prevent over activation of the T cell. Activated T cells also change their cell surface glycosylation profile.
The T cell receptor exists as a complex of several proteins. The actual T cell receptor is composed of two separate peptide chains, which are produced from the independent T cell receptor alpha and beta (TCRα and TCRβ) genes. The other proteins in the complex are the CD3 proteins: CD3εγ and CD3εδ heterodimers and, most important, a CD3ζ homodimer, which has a total of six ITAM motifs. The ITAM motifs on the CD3ζ can be phosphorylated by Lck and in turn recruit ZAP-70. Lck and/or ZAP-70 can also phosphorylate the tyrosines on many other molecules, not least CD28, LAT and SLP-76, which allows the aggregation of signalling complexes around these proteins.
Phosphorylated LAT recruits SLP-76 to the membrane, where it can then bring in PLC-γ, VAV1, Itk and potentially PI3K. PLC-γ cleaves PI(4,5)P2 on the inner leaflet of the membrane to create the active intermediaries diacylglycerol (DAG), inositol-1,4,5-trisphosphate (IP3); PI3K also acts on PIP2, phosphorylating it to produce phosphatidlyinositol-3,4,5-trisphosphate (PIP3). DAG binds and activates some PKCs. Most important in T cells is PKC-θ, critical for activating the transcription factors NF-κB and AP-1. IP3 is released from the membrane by PLC-γ and diffuses rapidly to activate calcium channel receptors on the ER, which induces the release of calcium into the cytosol. Low calcium in the endoplasmic reticulum causes STIM1 clustering on the ER membrane and leads to activation of cell membrane CRAC channels that allows additional calcium to flow into the cytosol from the extracellular space. This aggregated cytosolic calcium binds calmodulin, which can then activate calcineurin. Calcineurin, in turn, activates NFAT, which then translocates to the nucleus. NFAT is a transcription factor that activates the transcription of a pleiotropic set of genes, most notable, IL-2, a cytokine that promotes long-term proliferation of activated T cells.
PLC-γ can also initiate the NF-κB pathway. DAG activates PKC-θ, which then phosphorylates CARMA1, causing it to unfold and function as a scaffold. The cytosolic domains bind an adapter BCL10 via CARD (Caspase activation and recruitment domains) domains; that then binds TRAF6, which is ubiquitinated at K63. This form of ubiquitination does not lead to degradation of target proteins. Rather, it serves to recruit NEMO, IKKα and -β, and TAB1-2/ TAK1. TAK 1 phosphorylates IKK-β, which then phosphorylates IκB allowing for K48 ubiquitination: leads to proteasomal degradation. Rel A and p50 can then enter the nucleus and bind the NF-κB response element. This coupled with NFAT signaling allows for complete activation of the IL-2 gene.
While in most cases activation is dependent on TCR recognition of antigen, alternative pathways for activation have been described. For example, cytotoxic T cells have been shown to become activated when targeted by other CD8 T cells leading to tolerization of the latter.
In spring 2014, the T-Cell Activation in Space (TCAS) experiment was launched to the International Space Station on the SpaceX CRS-3 mission to study how "deficiencies in the human immune system are affected by a microgravity environment".
T cell activation is modulated by reactive oxygen species.
Antigen discrimination
A unique feature of T cells is their ability to discriminate between healthy and abnormal (e.g. infected or cancerous) cells in the body. Healthy cells typically express a large number of self derived pMHC on their cell surface and although the T cell antigen receptor can interact with at least a subset of these self pMHC, the T cell generally ignores these healthy cells. However, when these very same cells contain even minute quantities of pathogen derived pMHC, T cells are able to become activated and initiate immune responses. The ability of T cells to ignore healthy cells but respond when these same cells contain pathogen (or cancer) derived pMHC is known as antigen discrimination. The molecular mechanisms that underlie this process are controversial.
Clinical significance
Deficiency
Causes of T cell deficiency include lymphocytopenia of T cells and/or defects on function of individual T cells. Complete insufficiency of T cell function can result from hereditary conditions such as severe combined immunodeficiency (SCID), Omenn syndrome, and cartilage–hair hypoplasia. Causes of partial insufficiencies of T cell function include acquired immune deficiency syndrome (AIDS), and hereditary conditions such as DiGeorge syndrome (DGS), chromosomal breakage syndromes (CBSs), and B cell and T cell combined disorders such as ataxia-telangiectasia (AT) and Wiskott–Aldrich syndrome (WAS).
The main pathogens of concern in T cell deficiencies are intracellular pathogens, including Herpes simplex virus, Mycobacterium and Listeria. Also, fungal infections are also more common and severe in T cell deficiencies.
Cancer
Cancer of T cells is termed T-cell lymphoma, and accounts for perhaps one in ten cases of non-Hodgkin lymphoma. The main forms of T cell lymphoma are:
Extranodal T cell lymphoma
Cutaneous T cell lymphomas: Sézary syndrome and Mycosis fungoides
Anaplastic large cell lymphoma
Angioimmunoblastic T cell lymphoma
Exhaustion
T cell exhaustion is a poorly defined or ambiguous term. There are three approaches to its definition. "The first approach primarily defines as exhausted the cells that present the same cellular dysfunction (typically, the absence of an expected effector response). The second approach primarily defines as exhausted the cells that are produced by a given cause (typically, but not necessarily, chronic exposure to an antigen). Finally, the third approach primarily defines as exhausted the cells that present the same molecular markers (typically, programmed cell death protein 1 [PD-1])."
Dysfunctional T cells are characterized by progressive loss of function, changes in transcriptional profiles and sustained expression of inhibitory receptors. At first, cells lose their ability to produce IL-2 and TNFα, which is followed by the loss of high proliferative capacity and cytotoxic potential, and eventually leads to their deletion. Exhausted T cells typically indicate higher levels of CD43, CD69 and inhibitory receptors combined with lower expression of CD62L and CD127. Exhaustion can develop during chronic infections, sepsis and cancer. Exhausted T cells preserve their functional exhaustion even after repeated antigen exposure.
During chronic infection and sepsis
T cell exhaustion can be triggered by several factors like persistent antigen exposure and lack of CD4 T cell help. Antigen exposure also has effect on the course of exhaustion because longer exposure time and higher viral load increases the severity of T cell exhaustion. At least 2–4 weeks exposure is needed to establish exhaustion. Another factor able to induce exhaustion are inhibitory receptors including programmed cell death protein 1 (PD1), CTLA-4, T cell membrane protein-3 (TIM3), and lymphocyte activation gene 3 protein (LAG3). Soluble molecules such as cytokines IL-10 or TGF-β are also able to trigger exhaustion. Last known factors that can play a role in T cell exhaustion are regulatory cells. Treg cells can be a source of IL-10 and TGF-β and therefore they can play a role in T cell exhaustion. Furthermore, T cell exhaustion is reverted after depletion of Treg cells and blockade of PD1. T cell exhaustion can also occur during sepsis as a result of cytokine storm. Later after the initial septic encounter anti-inflammatory cytokines and pro-apoptotic proteins take over to protect the body from damage. Sepsis also carries high antigen load and inflammation. In this stage of sepsis T cell exhaustion increases. Currently there are studies aiming to utilize inhibitory receptor blockades in treatment of sepsis.
During transplantation
While during infection T cell exhaustion can develop following persistent antigen exposure after graft transplant similar situation arises with alloantigen presence. It was shown that T cell response diminishes over time after kidney transplant. These data suggest T cell exhaustion plays an important role in tolerance of a graft mainly by depletion of alloreactive CD8 T cells. Several studies showed positive effect of chronic infection on graft acceptance and its long-term survival mediated partly by T cell exhaustion. It was also shown that recipient T cell exhaustion provides sufficient conditions for NK cell transfer. While there are data showing that induction of T cell exhaustion can be beneficial for transplantation it also carries disadvantages among which can be counted increased number of infections and the risk of tumor development.
During cancer
During cancer T cell exhaustion plays a role in tumor protection. According to research some cancer-associated cells as well as tumor cells themselves can actively induce T cell exhaustion at the site of tumor. T cell exhaustion can also play a role in cancer relapses as was shown on leukemia. Some studies have suggested that it is possible to predict relapse of leukemia based on expression of inhibitory receptors PD-1 and TIM-3 by T cells. Many experiments and clinical trials have focused on immune checkpoint blockers in cancer therapy, with some of these approved as valid therapies that are now in clinical use. Inhibitory receptors targeted by those medical procedures are vital in T cell exhaustion and blocking them can reverse these changes.
| Biology and health sciences | Immune system | Biology |
170567 | https://en.wikipedia.org/wiki/Toxicity | Toxicity | Toxicity is the degree to which a chemical substance or a particular mixture of substances can damage an organism. Toxicity can refer to the effect on a whole organism, such as an animal, bacterium, or plant, as well as the effect on a substructure of the organism, such as a cell (cytotoxicity) or an organ such as the liver (hepatotoxicity). Sometimes the word is more or less synonymous with poisoning in everyday usage.
A central concept of toxicology is that the effects of a toxicant are dose-dependent; even water can lead to water intoxication when taken in too high a dose, whereas for even a very toxic substance such as snake venom there is a dose below which there is no detectable toxic effect. Toxicity is species-specific, making cross-species analysis problematic. Newer paradigms and metrics are evolving to bypass animal testing, while maintaining the concept of toxicity endpoints.
Etymology
In Ancient Greek medical literature, the adjective τοξικόν (meaning "toxic") was used to describe substances which had the ability of "causing death or serious debilitation or exhibiting symptoms of infection." The word draws its origins from the Greek noun τόξον (meaning "arc"), in reference to the use of bows and poisoned arrows as weapons.
English-speaking American culture has adopted several figurative usages for toxicity, often when describing harmful inter-personal relationships or character traits (e.g. "toxic masculinity").
History
Humans have a deeply rooted history of not only being aware of toxicity, but also taking advantage of it as a tool. Archaeologists studying bone arrows from caves of Southern Africa have noted the likelihood that some aging 72,000 to 80,000 years old were dipped in specially prepared poisons to increase their lethality. Although scientific instrumentation limitations make it difficult to prove concretely, archaeologists hypothesize the practice of making poison arrows was widespread in cultures as early as the paleolithic era. The San people of Southern Africa have managed to preserved this practice into the modern era, with the knowledge base to form complex mixtures from poisonous beetles and plant derived extracts, yielding an arrow-tip product with a shelf life beyond several months to a year.
Types
There are generally five types of toxicities: chemical, biological, physical, radioactive and behavioural.
Disease-causing microorganisms and parasites are toxic in a broad sense but are generally called pathogens rather than toxicants. The biological toxicity of pathogens can be difficult to measure because the threshold dose may be a single organism. Theoretically one virus, bacterium or worm can reproduce to cause a serious infection. If a host has an intact immune system, the inherent toxicity of the organism is balanced by the host's response; the effective toxicity is then a combination. In some cases, e.g. cholera toxin, the disease is chiefly caused by a nonliving substance secreted by the organism, rather than the organism itself. Such nonliving biological toxicants are generally called toxins if produced by a microorganism, plant, or fungus, and venoms if produced by an animal.
Physical toxicants are substances that, due to their physical nature, interfere with biological processes. Examples include coal dust, asbestos fibres or finely divided silicon dioxide, all of which can ultimately be fatal if inhaled. Corrosive chemicals possess physical toxicity because they destroy tissues, but are not directly poisonous unless they interfere directly with biological activity. Water can act as a physical toxicant if taken in extremely high doses because the concentration of vital ions decreases dramatically with too much water in the body. Asphyxiant gases can be considered physical toxicants because they act by displacing oxygen in the environment but they are inert, not chemically toxic gases.
Radiation can have a toxic effect on organisms.
Behavioral toxicity refers to the undesirable effects of essentially therapeutic levels of medication clinically indicated for a given disorder (DiMascio, Soltys and Shader, 1970). These undesirable effects include anticholinergic effects, alpha-adrenergic blockade, and dopaminergic effects, among others.
Measuring
Toxicity can be measured by its effects on the target (organism, organ, tissue or cell). Because individuals typically have different levels of response to the same dose of a toxic substance, a population-level measure of toxicity is often used which relates the probabilities of an outcome for a given individual in a population. One such measure is the . When such data does not exist, estimates are made by comparison to known similar toxic things, or to similar exposures in similar organisms. Then, "safety factors" are added to account for uncertainties in data and evaluation processes. For example, if a dose of a toxic substance is safe for a laboratory rat, one might assume that one-tenth that dose would be safe for a human, allowing a safety factor of 10 to allow for interspecies differences between two mammals; if the data are from fish, one might use a factor of 100 to account for the greater difference between two chordate classes (fish and mammals). Similarly, an extra protection factor may be used for individuals believed to be more susceptible to toxic effects such as in pregnancy or with certain diseases. Or, a newly synthesized and previously unstudied chemical that is believed to be very similar in effect to another compound could be assigned an additional protection factor of 10 to account for possible differences in effects that are probably much smaller. This approach is very approximate, but such protection factors are deliberately very conservative, and the method has been found to be useful in a wide variety of applications.
Assessing all aspects of the toxicity of cancer-causing agents involves additional issues, since it is not certain if there is a minimal effective dose for carcinogens, or whether the risk is just too small to see. In addition, it is possible that a single cell transformed into a cancer cell is all it takes to develop the full effect (the "one hit" theory).
It is more difficult to determine the toxicity of chemical mixtures than a pure chemical because each component displays its own toxicity, and components may interact to produce enhanced or diminished effects. Common mixtures include gasoline, cigarette smoke, and industrial waste. Even more complex are situations with more than one type of toxic entity, such as the discharge from a malfunctioning sewage treatment plant, with both chemical and biological agents.
The preclinical toxicity testing on various biological systems reveals the species-, organ- and dose-specific toxic effects of an investigational product. The toxicity of substances can be observed by (a) studying the accidental exposures to a substance (b) in vitro studies using cells/ cell lines (c) in vivo exposure on experimental animals. Toxicity tests are mostly used to examine specific adverse events or specific endpoints such as cancer, cardiotoxicity, and skin/eye irritation. Toxicity testing also helps calculate the No Observed Adverse Effect Level (NOAEL) dose and is helpful for clinical studies.
Classification
For substances to be regulated and handled appropriately they must be properly classified and labelled. Classification is determined by approved testing measures or calculations and has determined cut-off levels set by governments and scientists (for example, no-observed-adverse-effect levels, threshold limit values, and tolerable daily intake levels). Pesticides provide the example of well-established toxicity class systems and toxicity labels. While currently many countries have different regulations regarding the types of tests, numbers of tests and cut-off levels, the implementation of the Globally Harmonized System has begun unifying these countries.
Global classification looks at three areas: Physical Hazards (explosions and pyrotechnics), Health Hazards and environmental hazards.
Health hazards
The types of toxicities where substances may cause lethality to the entire body, lethality to specific organs, major/minor damage, or cause cancer. These are globally accepted definitions of what toxicity is. Anything falling outside of the definition cannot be classified as that type of toxicant.
Acute toxicity
Acute toxicity looks at lethal effects following oral, dermal or inhalation exposure. It is split into five categories of severity where Category 1 requires the least amount of exposure to be lethal and Category 5 requires the most exposure to be lethal. The table below shows the upper limits for each category.
Note: The undefined values are expected to be roughly equivalent to the category 5 values for oral and dermal administration.
Other methods of exposure and severity
Skin corrosion and irritation are determined through a skin patch test analysis, similar to an allergic inflammation patch test. This examines the severity of the damage done; when it is incurred and how long it remains; whether it is reversible and how many test subjects were affected.
Skin corrosion from a substance must penetrate through the epidermis into the dermis within four hours of application and must not reverse the damage within 14 days. Skin irritation shows damage less severe than corrosion if: the damage occurs within 72 hours of application; or for three consecutive days after application within a 14-day period; or causes inflammation which lasts for 14 days in two test subjects. Mild skin irritation is minor damage (less severe than irritation) within 72 hours of application or for three consecutive days after application.
Serious eye damage involves tissue damage or degradation of vision which does not fully reverse in 21 days. Eye irritation involves changes to the eye which do fully reverse within 21 days.
Other categories
Respiratory sensitizers cause breathing hypersensitivity when the substance is inhaled.
A substance which is a skin sensitizer causes an allergic response from a dermal application.
Carcinogens induce cancer, or increase the likelihood of cancer occurring.
Neurotoxicity is a form of toxicity in which a biological, chemical, or physical agent produces an adverse effect on the structure or function of the central or peripheral nervous system. It occurs when exposure to a substance – specifically, a neurotoxin or neurotoxicant– alters the normal activity of the nervous system in such a way as to cause permanent or reversible damage to nervous tissue.
Reproductively toxic substances cause adverse effects in either sexual function or fertility to either a parent or the offspring.
Specific-target organ toxins damage only specific organs.
Aspiration hazards are solids or liquids which can cause damage through inhalation.
Environmental hazards
An Environmental hazard can be defined as any condition, process, or state adversely affecting the environment. These hazards can be physical or chemical, and present in air, water, and/or soil. These conditions can cause extensive harm to humans and other organisms within an ecosystem.
Common types of environmental hazards
Water: detergents, fertilizer, raw sewage, prescription medication, pesticides, herbicides, heavy metals, PCBs
Soil: heavy metals, herbicides, pesticides, PCBs
Air: particulate matter, carbon monoxide, sulfur dioxide, nitrogen dioxide, asbestos, ground-level ozone, lead (from aircraft fuel, mining, and industrial processes)
The EPA maintains a list of priority pollutants for testing and regulation.
Occupational hazards
Workers in various occupations may be at a greater level of risk for several types of toxicity, including neurotoxicity. The expression "Mad as a hatter" and the "Mad Hatter" of the book Alice in Wonderland derive from the known occupational toxicity of hatters who used a toxic chemical for controlling the shape of hats. Exposure to chemicals in the workplace environment may be required for evaluation by industrial hygiene professionals.
Hazards for small businesses
Hazards from medical waste and prescription disposal
Hazards in the arts
Hazards in the arts have been an issue for artists for centuries, even though the toxicity of their tools, methods, and materials was not always adequately realized. Lead and cadmium, among other toxic elements, were often incorporated into the names of artist's oil paints and pigments, for example, "lead white" and "cadmium red".
20th-century printmakers and other artists began to be aware of the toxic substances, toxic techniques, and toxic fumes in glues, painting mediums, pigments, and solvents, many of which in their labelling gave no indication of their toxicity. An example was the use of xylol for cleaning silk screens. Painters began to notice the dangers of breathing painting mediums and thinners such as turpentine. Aware of toxicants in studios and workshops, in 1998 printmaker Keith Howard published Non-Toxic Intaglio Printmaking which detailed twelve innovative Intaglio-type printmaking techniques including photo etching, digital imaging, acrylic-resist hand-etching methods, and introducing a new method of non-toxic lithography.
Mapping environmental hazards
There are many environmental health mapping tools. TOXMAP is a Geographic Information System (GIS) from the Division of Specialized Information Services of the United States National Library of Medicine (NLM) that uses maps of the United States to help users visually explore data from the United States Environmental Protection Agency's (EPA) Toxics Release Inventory and Superfund programs. TOXMAP is a resource funded by the US Federal Government. TOXMAP's chemical and environmental health information is taken from NLM's Toxicology Data Network
(TOXNET) and PubMed, and from other authoritative sources.
Aquatic toxicity
Aquatic toxicity testing subjects key indicator species of fish or crustacea to certain concentrations of a substance in their environment to determine the lethality level. Fish are exposed for 96 hours while crustacea are exposed for 48 hours. While GHS does not define toxicity past 100 mg/L, the EPA currently lists aquatic toxicity as "practically non-toxic" in concentrations greater than 100 ppm.
Note: A category 4 is established for chronic exposure, but simply contains any toxic substance which is mostly insoluble, or has no data for acute toxicity.
Factors influencing toxicity
Toxicity of a substance can be affected by many different factors, such as the pathway of administration (whether the toxicant is applied to the skin, ingested, inhaled, injected), the time of exposure (a brief encounter or long term), the number of exposures (a single dose or multiple doses over time), the physical form of the toxicant (solid, liquid, gas), the concentration of the substance, and in the case of gases, the partial pressure (at high ambient pressure, partial pressure will increase for a given concentration as a gas fraction), the genetic makeup of an individual, an individual's overall health, and many others. Several of the terms used to describe these factors have been included here.
Acute exposure A single exposure to a toxic substance which may result in severe biological harm or death; acute exposures are usually characterized as lasting no longer than a day.
Chronic exposure Continuous exposure to a toxicant over an extended period of time, often measured in months or years; it can cause irreversible side effects.
Alternatives to dose-response framework
Considering the limitations of the dose-response concept, a novel Abstract Drug Toxicity Index (DTI) has been proposed recently. DTI redefines drug toxicity, identifies hepatotoxic drugs, gives mechanistic insights, predicts clinical outcomes and has potential as a screening tool.
| Biology and health sciences | Miscellaneous | null |
170578 | https://en.wikipedia.org/wiki/Vermicompost | Vermicompost | Vermicompost (vermi-compost) is the product of the decomposition process using various species of worms, usually red wigglers, white worms, and other earthworms, to create a mixture of decomposing vegetable or food waste, bedding materials, and vermicast. This process is called vermicomposting, with the rearing of worms for this purpose is called vermiculture.
Vermicast (also called worm castings, worm humus, worm poop, worm manure, or worm faeces) is the end-product of the breakdown of organic matter by earthworms. These excreta have been shown to contain reduced levels of contaminants and a higher saturation of nutrients than the organic materials before vermicomposting.
Vermicompost contains water-soluble nutrients which may be extracted as vermiwash and is an excellent, nutrient-rich organic fertilizer and soil conditioner. It is used in gardening and sustainable, organic farming.
Vermicomposting can also be applied for treatment of sewage. A variation of the process is vermifiltration (or vermidigestion) which is used to remove organic matter, pathogens, and oxygen demand from wastewater or directly from blackwater of flush toilets.
Overview
Vermicomposting has gained popularity in both industrial and domestic settings because, as compared with conventional composting, it provides a way to treat organic wastes more quickly. In manure composing, the use of vermicomposting generates products that have lower salinity levels, as well as a more neutral pH.
The earthworm species (or composting worms) most often used are red wigglers (Eisenia fetida or Eisenia andrei), though European nightcrawlers (Eisenia hortensis, synonym Dendrobaena veneta) and red earthworm (Lumbricus rubellus) could also be used. Red wigglers are recommended by most vermicomposting experts, as they have some of the best appetites and breed very quickly. Users refer to European nightcrawlers by a variety of other names, including dendrobaenas, dendras, Dutch nightcrawlers, and Belgian nightcrawlers.
Containing water-soluble nutrients, vermicompost is a nutrient-rich organic fertilizer and soil conditioner in a form that is relatively easy for plants to absorb. Worm castings are sometimes used as an organic fertilizer. Because the earthworms grind and uniformly mix minerals in simple forms, plants need only minimal effort to obtain them. The worms' digestive systems create environments that allow certain species of microbes to thrive to help create a "living" soil environment for plants. The fraction of soil which has gone through the digestive tract of earthworms is called the drilosphere.
Vermicomposting is a common practice in permaculture.
Vermiwash can also be obtained from the liquid potion of vermicompost. Vermiwash is found to contain enzyme cocktail of proteases, amylases, urease and phosphatase. Microbiological study of vermiwash reveals that it contains nitrogen-fixing bacteria like Azotobactrer sp., Agrobacterium sp. and Rhizobium sp. and some phosphate solublizing bacteria. Laboratory scale trial shows effectiveness of vermiwash on plant growth.
Design considerations
Suitable worm species
All worms make compost but some species are not suitable for this purpose. Vermicompost worms are generally epigean. Species most often used for composting include:
Eisenia fetida (Europe), the red wiggler or tiger worm. Closely related to Eisenia andrei, which is also usable.
Eisenia hortensis (Europe), European nightcrawlers, prefers high C:N material.
Eudrilus eugeniae (West Africa), African Nightcrawlers. Useful in the tropics.
Perionyx excavatus (South and East Asia), blueworms. May be used in the tropics and subtropics.
Lampito mauritii (Southern Asia), used locally.
These species commonly are found in organic-rich soils throughout Europe and North America and live in rotting vegetation, compost, and manure piles. As they are shallow-dwelling and feed on decomposing plant matter in the soil, they adapt easily to live on food or plant waste in the confines of a worm bin. Some species are considered invasive in some areas, so they should be avoided (see earthworms as invasive species for a list).
Composting worms are available to order online, from nursery mail-order suppliers or angling shops where they are sold as bait. They can also be collected from compost and manure piles. These species are not the same worms that are found in ordinary soil or on pavement when the soil is flooded by water.
The following species are not recommended:
Lumbricus rubellus and Lumbricus terrestris (Europe). The two closely related species are anecic: they like to burrow underground and come up for food. As a result, they adapt poorly to shallow compost bins and should be avoided. They are also invasive in North America.
Large scale
Large-scale vermicomposting is practiced in Canada, Italy, Japan, India, Malaysia, the Philippines, and the United States. The vermicompost may be used for farming, landscaping, to create compost tea, or for sale. Some of these operations produce worms for bait and/or home vermicomposting.
There are two main methods of large-scale vermicomposting, windrow and raised bed. Some systems use a windrow, which consists of bedding materials for the earthworms to live in and acts as a large bin; organic material is added to it. Although the windrow has no physical barriers to prevent worms from escaping, in theory they should not, due to an abundance of organic matter for them to feed on. Often windrows are used on a concrete surface to prevent predators from gaining access to the worm population.
The windrow method and compost windrow turners were developed by Fletcher Sims Jr. of the Compost Corporation in Canyon, Texas. The Windrow Composting system is noted as a sustainable, cost-efficient way for farmers to manage dairy waste.
The second type of large-scale vermicomposting system is the raised bed or flow-through system. Here the worms are fed an inch of "worm chow" across the top of the bed, and an inch of castings are harvested from below by pulling a breaker bar across the large mesh screen which forms the base of the bed.
Because red worms are surface dwellers constantly moving towards the new food source, the flow-through system eliminates the need to separate worms from the castings before packaging. Flow-through systems are well suited to indoor facilities, making them the preferred choice for operations in colder climates.
Small scale
For vermicomposting at home, a large variety of bins are commercially available, or a variety of adapted containers may be used. They may be made of old plastic containers, wood, Styrofoam, or metal containers. The design of a small bin usually depends on where an individual wishes to store the bin and how they wish to feed the worms.
Some materials are less desirable than others in worm bin construction. Metal containers often conduct heat too readily, are prone to rusting, and may release heavy metals into the vermicompost. Styrofoam containers may release chemicals into the organic material. Some cedars, yellow cedar, and redwood contain resinous oils that may harm worms, although western red cedar has excellent longevity in composting conditions. Hemlock is another inexpensive and fairly rot-resistant wood species that may be used to build worm bins.
Bins need holes or mesh for aeration. Some people add a spout or holes in the bottom for excess liquid to drain into a tray for collection. The most common materials used are plastic: recycled polyethylene and polypropylene and wood. Worm compost bins made from plastic are ideal, but require more drainage than wooden ones because they are non-absorbent. However, wooden bins will eventually decay and need to be replaced.
Small-scale vermicomposting is well-suited to turn kitchen waste into high-quality soil amendments, where space is limited. Worms can decompose organic matter without the additional human physical effort (turning the bin) that bin composting requires.
Composting worms which are detritivorous (eaters of trash), such as the red wiggler Eisenia fetida, are epigeic (surface dwellers) and together with symbiotic associated microbes are the ideal vectors for decomposing food waste. Common earthworms such as Lumbricus terrestris are anecic (deep burrowing) species and hence unsuitable for use in a closed system. Other soil species that contribute include insects, other worms and molds.
Climate and temperature
There may be differences in vermicomposting method depending on the climate. It is necessary to monitor the temperatures of large-scale bin systems (which can have high heat-retentive properties), as the raw materials or feedstocks used can compost, heating up the worm bins as they decay and killing the worms.
The most common worms used in composting systems, redworms (Eisenia fetida, Eisenia andrei, and Lumbricus rubellus) feed most rapidly at temperatures of . They can survive at . Temperatures above may harm them. This temperature range means that indoor vermicomposting with redworms is possible in all but tropical climates. Other worms like Perionyx excavatus are suitable for warmer climates. If a worm bin is kept outside, it should be placed in a sheltered position away from direct sunlight and insulated against frost in winter.
Feedstock
There are few food wastes that vermicomposting cannot compost, although meat waste and dairy products are likely to putrefy, and in outdoor bins can attract vermin. Green waste should be added in moderation to avoid heating the bin.
Small-scale or home systems
Such systems usually use kitchen and garden waste, using "earthworms and other microorganisms to digest organic wastes, such as kitchen scraps".
This includes:
All fruits and vegetables (including citrus, in limited quantities)
Vegetable and fruit peels and ends
Coffee grounds and filters
Tea bags (even those with high tannin levels)
Grains such as bread, cracker and cereal (including moldy and stale)
Eggshells (rinsed off)
Leaves and grass clippings (not sprayed with pesticides)
Newspapers (most inks used in newspapers are not toxic)
Paper toweling (which has not been used with cleaners or chemicals)
Large-scale or commercial
Such vermicomposting systems need reliable sources of large quantities of food.
Systems presently operating use:
Dairy cow or pig manure
Sewage sludge
Brewery waste
Cotton mill waste
Agricultural waste
Food processing and grocery waste
Cafeteria waste
Grass clippings and wood chips
Harvesting
Factors affecting the speed of composting include the climate and the method of composting. There are signs to look for to determine whether compost is finished. The finished compost would have an ambient temperature, dark color, and be as moist as a damp sponge. Towards the end of the process, bacteria slow down the rate of metabolizing food or stop completely. There is the possibility of some solid organic matter still being present in the compost at this point, but it could stay in and continue decomposing for the next couple of years unless removed. The compost should be allowed to cure after finished to allow acids to be removed over time so it becomes more neutral, which could take up to three months and results in the compost being more consistent in size. Elevating the maturing compost off the ground can prevent unwanted plant growth. It compost should consistently be slightly damp and should be aerated but does not need to be turned. The curing process can be done in a storage bin or on a tarp.
Methods
Vermicompost is ready for harvest when it contains few-to-no scraps of uneaten food or bedding. There are several methods of harvesting from small-scale systems: "dump and hand sort", "let the worms do the sorting", "alternate containers" and "divide and dump." These differ on the amount of time and labor involved and whether the vermicomposter wants to save as many worms as possible from being trapped in the harvested compost.
The pyramid method of harvesting worm compost is commonly used in small-scale vermicomposting, and is considered the simplest method for single layer bins. In this process, compost is separated into large clumps, which is placed back into composting for further breakdown, and lighter compost, with which the rest of the process continues. This lighter mix is placed into small piles on a tarp under the sunlight. The worms instinctively burrow to the bottom of the pile. After a few minutes, the top of the pyramid is removed repeatedly, until the worms are again visible. This repeats until the mound is composed mostly of worms.
When harvesting the compost, it is possible to separate eggs and cocoons and return them to the bin, thereby ensuring new worms are hatched. Cocoons are small, lemon-shaped yellowish objects that can usually be seen with the naked eye. The cocoons can hold up to 20 worms (though 2–3 is most common). Cocoons can lay dormant for as long as two years if conditions are not conducive for hatching.
Properties
Vermicompost has been shown to be richer in many nutrients than compost produced by other composting methods. It has also outperformed a commercial plant medium with nutrients added, but levels of magnesium required adjustment, as did pH.
However, in one study it has been found that homemade backyard vermicompost was lower in microbial biomass, soil microbial activity, and yield of a species of ryegrass than municipal compost.
It is rich in microbial life which converts nutrients already present in the soil into plant-available forms.
Unlike other compost, worm castings also contain worm mucus which helps prevent nutrients from washing away with the first watering and holds moisture better than plain soil.
Increases in the total nitrogen content in vermicompost, an increase in available nitrogen and phosphorus, a decrease in potassium, as well as the increased removal of heavy metals from sludge and soil have been reported. The reduction in the bioavailability of heavy metals has been observed in a number of studies.
Benefits of vermicomposting
Soil
Improves soil aeration
Enriches soil with micro-organisms (adding enzymes such as phosphatase and cellulase)
Microbial activity in worm castings is 10 to 20 times higher than in the soil and organic matter that the worm ingests
Attracts deep-burrowing earthworms already present in the soil
Improves water holding capacity
Plant growth
Enhances germination, plant growth and crop yield
It helps in root and plant growth
Enriches soil organisms (adding plant hormones such as auxins and gibberellic acid)
Economic
Biowastes conversion reduces waste flow to landfills
Elimination of biowastes from the waste stream reduces contamination of other recyclables collected in a single bin (a common problem in communities practicing single-stream recycling)
Creates low-skill jobs at local level
Low capital investment and relatively simple technologies make vermicomposting practical for less-developed agricultural regions
Environmental
Helps to close the "metabolic gap" through recycling waste on-site
Large systems often use temperature control and mechanized harvesting, however other equipment is relatively simple and does not wear out quickly
Production reduces greenhouse gas emissions such as methane and nitric oxide (produced in landfills or incinerators when not composted).
Uses
Soil conditioner
Vermicompost can be mixed directly into the soil, or mixed with water to make a liquid fertilizer known as worm tea.
The light brown waste liquid, or leachate, that drains into the bottom of some vermicomposting systems is not to be confused with worm tea. It is an uncomposted byproduct from when water-rich foods break down and may contain pathogens and toxins. It is best discarded or applied back to the bin when added moisture is needed for further processing.
The pH, nutrient, and microbial content of these fertilizers varies upon the inputs fed to worms. Pulverized limestone, or calcium carbonate can be added to the system to raise the pH.
Operation and maintenance
Smells
When closed, a well-maintained bin is odorless; when opened, it should have little smell—if any smell is present, it is earthy. The smell may also depend on the type of composted material added to the bin. An unhealthy worm bin may smell, potentially due to low oxygen conditions. Worms require gaseous oxygen. Oxygen can be provided by airholes in the bin, occasional stirring of bin contents, and removal of some bin contents if they become too deep or too wet. If decomposition becomes anaerobic from excess wet feedstock added to the bin, or the layers of food waste have become too deep, the bin will begin to smell of ammonia.
Moisture
Moisture must be maintained above 50%, as lower moisture content will not support worm respiration and can increase worm mortality. Operating moisture-content range should be between 70 and 90%, with a suggested content of 70–80% for vermicomposting operations. If decomposition has become anaerobic, to restore healthy conditions and prevent the worms from dying, excess waste water must be reduced and the bin returned to a normal moisture level. To do this, first reduce addition of food scraps with a high moisture content and second, add fresh, dry bedding such as shredded newspaper to your bin, mixing it in well.
Pest species
Pests such as rodents and flies are attracted by certain materials and odors, usually from large amounts of kitchen waste, particularly meat. Eliminating the use of meat or dairy product in a worm bin decreases the possibility of pests.
Predatory ants can be a problem in African countries.
In warm weather, fruit and vinegar flies breed in the bins if fruit and vegetable waste is not thoroughly covered with bedding. This problem can be avoided by thoroughly covering the waste by at least of bedding. Maintaining the correct pH (close to neutral) and water content of the bin (just enough water where squeezed bedding drips a couple of drops) can help avoid these pests as well.
Worms escaping
Worms generally stay in the bin, but may try to leave the bin when first introduced, or often after a rainstorm when the humidity outside is high. Maintaining adequate conditions in the worm bin and putting a light over the bin when first introducing worms should eliminate this problem.
Nutrient levels
Commercial vermicomposters test and may amend their products to produce consistent quality and results. Because the small-scale and home systems use a varied mix of feedstocks, the nitrogen, phosphorus, and potassium (NPK) content of the resulting vermicompost will also be inconsistent. NPK testing may be helpful before the vermicompost or tea is applied to the garden.
In order to avoid over-fertilization issues, such as nitrogen burn, vermicompost can be diluted as a tea 50:50 with water, or as a solid can be mixed in 50:50 with potting soil.
Additionally, the mucous layer created by worms which surrounds their castings allows for a "time release" effect, meaning not all nutrients are released at once. This also reduces the risk of burning the plants, as is common with the use and overuse of commercial fertilizers.
Application examples
Vermicomposting is widely used in North America for on-site institutional processing of food scraps, such as in hospitals, universities, shopping malls, and correctional facilities. Vermicomposting is used for medium-scale on-site institutional organic material recycling, such as for food scraps from universities and shopping malls. It is selected either as a more environmentally friendly choice than conventional disposal, or to reduce the cost of commercial waste removal.
From 20 July 2020, the State Government of Chhattisgarh India started buying cow dung under the "Godhan Nyay Yojana" Scheme. Cow dung procured under this scheme will be utilised for the production of vermicompost fertilizer.
| Technology | Agronomical techniques | null |
170660 | https://en.wikipedia.org/wiki/Angora%20goat | Angora goat | The Angora or Ankara is a Turkish breed of domesticated goat. It produces the lustrous fibre known as mohair. It is widespread in many countries of the world. Many breeds derive from it, among them the Indian Mohair, the Soviet Mohair, the Angora-Don of the Russian Federation and the Pygora in the United States.
History
The origin of the Angora is not known. The earliest Western description may be that published in 1555 by Pierre Belon, who while travelling from Heraclea to Konya in southern Turkey had seen goats with snow-white "... wool so delicate that one would judge it finer than silk ...".
Angora goats were depicted on the reverse of the Turkish 50 lira banknote from 1938 to 1952.
In 1960 there were over 6 million Angora goats in Turkey; the population subsequently dropped sharply. In 2004 the total goat population of the country was approximately 7.2 million; of these, just over 5% were of Angora stock, while the remainder were hair goats. A conservation programme for the Angora was established in 2003.
Characteristics
The Angora is a moderately small goat, standing about at the withers. It is slender, elegant and light-framed; the head is small, with semi-lop ears. It is usually horned; in billies the horns are commonly long, twisted and strong. With the exception of the face and legs, the animal is entirely covered in a coat of long ringlets of fine and lustrous mohair. This is not goat hair as seen on other breeds, but the down or undercoat which, in this breed only, grows much longer than the outer hair coat. The face and coat are normally white, but – particularly in southern Turkey – black, brown and grey animals also occur.
Use
The goats are reared either for mohair or for their goat's meat. Mohair is not as fine as cashmere, but yields are much higher. Unlike cashmere, which is obtained by combing the coat of the goat, mohair is obtained by shearing; this is commonly done twice per year. In 2010 approximately half of all mohair production was in South Africa; Argentina and Lesotho were also major producers, followed by the United States, Turkey, Australia and New Zealand.
In some other countries the Angora is reared for its meat, which is succulent and tender, and which in the early twentieth century was described as the best of its kind in the world.
| Biology and health sciences | Goats | Animals |
170690 | https://en.wikipedia.org/wiki/Cuniculture | Cuniculture | Cuniculture is the agricultural practice of breeding and raising domestic rabbits as livestock for their meat, fur, or wool. Cuniculture is also employed by rabbit fanciers and hobbyists in the development and betterment of rabbit breeds and the exhibition of those efforts. Scientists practice cuniculture in the use and management of rabbits as model organisms in research. Cuniculture has been practiced all over the world since at least the 5th century.
History
Early husbandry
An abundance of ancient rabbits may have played a part in the naming of Spain. Phoenician sailors visiting its coast around the 12th century BC mistook the European rabbit for the familiar rock hyrax (Procavia capensis) of their homeland. They named their discovery , meaning 'land [or island] of hyraxes'. A theory exists (though it is somewhat controversial) that a corruption of this name used by the Romans became Hispania, the Latin name for the Iberian Peninsula.
Domestication of the European rabbit rose slowly from a combination of game-keeping and animal husbandry. Among the numerous foodstuffs imported by sea to Rome during her domination of the Mediterranean were shipments of rabbits from Spain. Romans also imported ferrets for rabbit hunting, and the Romans then distributed rabbits and the habit of rabbit keeping to the rest of Italy, to France, and then across the Roman Empire, including the British Isles. Rabbits were kept in both walled areas as well as more extensively in game-preserves. In the British Isles, these preserves were known as warrens or , and rabbits were known as , to differentiate them from the similar hares. The term warren was also used as a name for the location where hares, partridges and pheasants were kept, under the watch of a game keeper called a warrener. In order to confine and protect the rabbits, a wall or thick hedge might be constructed around the warren, or a warren might be established on an island. A warrener was responsible for controlling poachers and other predators and would collect the rabbits with snares, nets, hounds (such as greyhounds), or by hunting with ferrets. With the rise of falconry, hawks and falcons were also used to collect rabbits and hares.
Domestication
While under the warren system, rabbits were managed and harvested, but not domesticated. The practice of rabbit domestication also came from Rome. Christian monasteries throughout Europe and the Middle East kept rabbits since at least the 5th century. While rabbits might be allowed to wander freely within the monastery walls, a more common method was the employment of rabbit courts or rabbit pits. A rabbit court was a walled area lined with brick and cement, while a pit was similar, although less well-lined and more sunken. Individual boxes or burrow-spaces could line the wall. Rabbits would be kept in a group in these pits or courts, and individuals collected when desired for eating or pelts. Rabbit keepers transferred rabbits to individual hutches or pens for easy cleaning, handling, or for selective breeding, as pits did not allow keepers to perform these tasks. Hutches or pens were originally made of wood, but are now more frequently made of metal in order to allow for better sanitation.
Early breeds
Rabbits were typically kept as part of the household livestock by peasants and villagers throughout Europe. Husbandry of the rabbits, including collecting weeds and grasses for fodder, typically fell to the children of the household or farmstead. These rabbits were largely 'common' or 'meat' rabbits and not of a particular breed, although regional strains and types did arise. Some of these strains remain as regional breeds, such as the Gotland of Sweden, while others, such as the Land Kaninchen, a spotted rabbit of Germany, have become extinct. Another rabbit type that standardized into a breed was the Brabancon, a meat rabbit of the region of Limbourg and what is now Belgium. Rabbits of this breed were bred for the Ostend port market, destined for London markets. The development of the refrigerated shipping vessels led to the eventual collapse of the European meat rabbit trade, as the over-populated feral rabbits in Australia could now be harvested and sold. The Brabancon is now considered extinct, although a descendant, the Dutch breed, remains a popular small rabbit for the pet trade.
In addition to being harvested for meat, properly prepared rabbit pelts were also an economic factor. Both wild rabbits and domestic rabbit pelts were valued, and it followed that pelts of particular rabbits would be more highly prized. As far back as 1631, price differentials were noted between ordinary rabbit pelts and the pelts of quality 'riche' rabbit in the Champagne region of France. (This regional type would go on to be recognized as the , the 'silver rabbit of Champagne'.)
Among the earliest of the commercial breeds was the Angora, which some say may have developed in the Carpathian Mountains. They made their way to England, where during the rule of King Henry VIII, laws banned the exportation of long-haired rabbits as they were a national treasure. In 1723, long haired rabbits were imported to southern France by English sailors, who described the animals as originally coming from the Angora region of Turkey. Thus two distinct strains arose, one in France and one in England.
Expansion around the globe
European explorers and sailors took rabbits with them to new ports around the world, and brought new varieties back to Europe and England with them. With the second voyage of Christopher Columbus in 1494, European domestic livestock were brought to the New World. Rabbits, along with goats and other hardy livestock, were frequently released on islands to produce a food supply for later ships. The importations occasionally met with disastrous results, such as in the devastation in Australia. While cattle and horses were used across the socio-economic spectrum, and especially were concentrated among the wealthy, rabbits were kept by lower-income classes and peasants. This is reflected in the names given to the breeds that eventually arose in the colonized areas. From the Santa Duromo mountains of Brazil comes the Rustico, which is known in the United States as the Brazilian rabbit. The Criollo rabbit comes from Mexico.
International commercial use
With the rise of scientific animal breeding in the late 1700s, led by Robert Bakewell (among others), distinct livestock breeds were developed for specific purposes.
Rabbits were among the last of the domestic animals to have these principles applied to them, but the rabbit's rapid reproductive cycle allowed for marked progress towards a breeding goal in a short period of time. Additionally, rabbits could be kept on a small area, with a single person caring for over 300 breeding does on an acre of land. Rabbit breeds were developed by individuals, cooperatives, and by national breeding centers. To meet various production goals, rabbits were exported around the world. One of the most notable import events was the introduction of the Belgian Hare breed of rabbit from Europe to the United States, beginning in 1888. This led to a short-lived "boom" in rabbit breeding, selling, and speculation, when a quality breeding animal could bring $75 to $200. (For comparison, the average daily wage at the time was approximately $1.) In 1900, a single animal-export company recorded 6,000 rabbits successfully shipped to the United States and Canada.
Science played another role in rabbit raising, this time with rabbits themselves as the tools used for scientific advancement. Beginning with Louis Pasteur's experiments in rabies in the later half of the nineteenth century, rabbits have been used as models to investigate various medical and biological problems, including the transmission of disease and protective antiserums. Production of quality animals for meat sale and scientific experimentation has driven a number of advancements in rabbit husbandry and nutrition. While early rabbit keepers were limited to local and seasonal foodstuffs, which did not permit the maximization of production, health or growth, by 1930 researchers were conducting experiments in rabbit nutrition, similar to the experiments that had isolated vitamins and other nutritional components. This eventually resulted in the development of various recipes for pelleted rabbit diets. Gradual refinement of diets has resulted in the widespread availability of pelleted feeds, which increase yield, reduce waste, and promote rabbit health, particularly maternal breeding health.
Rise of the fancy
The final leg of rabbit breeding—beyond meat, wool, fur, and laboratory use—was the breeding of 'fancy' animals as pets and curiosities. The term 'fancy' was originally applied to long-eared 'lop' rabbits, as they were the first type to be bred for exhibition.
Such rabbits were first admitted to agricultural shows in England in the 1820s, and in 1840 a club was formed for the promotion and regulation of exhibitions for "Fancy Rabbits". In 1918, a new group formed to promote the fur breeds, originally just the Beveren and Havana breeds.
This club eventually expanded to become the British Rabbit Council. Meanwhile, in the United States, clubs promoting various breeds were chartered in the 1880s, and the National Pet Stock Association was formed in 1910. This organization would become the American Rabbit Breeders Association. Thousands of rabbit shows take place each year and are sanctioned in Canada, Mexico, Malaysia, Indonesia and the United States by ARBA.
With the advent of national-level organizations, rabbit breeders had a framework for establishing breeds and varieties utilizing recognized standards, and breeding for rabbit exhibitions began to expand rapidly. Such organizations and associations were also established across Europe—most notably in Germany, France, and Scandinavia—allowing for the recognition of local breeds (many of which shared similar characteristics across national borders) and for the preservation of stock during disruptions such as World WarI and World WarII.
Closely overlapping with breeding for exhibition and for fur has been the breeding of rabbits for the pet trade. While rabbits have been kept as companions for centuries, the sale of rabbits as pets began to rise in the last half of the twentieth century. This may have been, in part, because rabbits require less physical space than dogs or cats, and do not require a specialized habitat like goldfish. Several breeds of rabbit—such as the Holland Lop, the Polish, the Netherland Dwarf, and the Lionhead—have been specifically bred for the pet trade. Traits common to many popular pet breeds are small size, "dwarf" (or neotenic) features, plush or fuzzy coats, and an array of coat colors and patterns.
Modern farming
Outside of the exhibition circles, rabbit raising remained a small-scale but persistent household and farm endeavor, in many locations unregulated by the rules that governed the production of larger livestock. With the ongoing urbanization of populations worldwide, rabbit raising gradually declined, but saw resurgences in both Europe and North America during World WarII, in conjunction with victory gardens. Eventually, farmers across Europe and in the United States began to approach cuniculture with the same scientific principles as had already been applied to the production of grains, poultry, and hoofed livestock. National agriculture breeding stations were established to improve local rabbit strains and to introduce more productive breeds. National breeding centers focused on developing strains for production purposes, including meat, pelts, and wool.
These gradually faded from prominence in the United States, but remained viable longer in Europe. Meanwhile, rabbit raising for local markets gained prominence in developing nations as an economical means of producing protein. Various aid agencies promote the use of rabbits as livestock. The animals are particularly useful in areas where women are limited in employment outside the household, because rabbits can be kept successfully in small areas. These same factors have contributed to the increased popularity of rabbits as "backyard livestock" among locavores and homesteaders in more developed countries in North America and Europe. The addition of rabbits to the watchlist of endangered heritage breeds that is kept by The Livestock Conservancy has also led to increased interest from livestock conservationists. In contrast, throughout Asia (and particularly in China) rabbits are increasingly being raised and sold for export around the world.
The World Rabbit Science Association (WRSA), formed in 1976, was established "to facilitate in all possible ways the exchange of knowledge and experience among persons in all parts of the world who are contributing to the advancement of the various branches of the rabbit industry". The WRSA organizes a world conference every four years.
Present day (2000–present)
Approximately 1.2 billion rabbits are slaughtered each year for meat worldwide. In more recent years and in some countries, cuniculture has come under pressure from animal rights activists on several fronts. The use of animals, including rabbits, in scientific experiments has been subject to increased scrutiny in developed countries. Increasing regulation has raised the cost of producing animals for this purpose, and made other experimental options more attractive. Other researchers have abandoned investigations which required animal models. Meanwhile, various rescue groups under the House Rabbit Society umbrella have taken an increasingly strident stance against any breeding of rabbits (even as food in developing countries) on the grounds that it contributes to the number of mistreated, unwanted or abandoned animals.
The growth of homesteaders and smallholders has led to the rise of visibility of rabbit raisers in geographic areas where they have not been previously present. This has led to zoning conflicts over the regulation of butchering and waste management. Conflicts have also arisen with House Rabbit Society organizations as well as ethical vegetarians and vegans concerning the use of rabbits as meat and fur animals rather than as pets. Conversely, many homesteaders cite concern with animal welfare in intensive farming of beef, pork and poultry as a significant factor in choosing to raise rabbits for meat.
In August 2022, an animal rights campaign group in the UK called "Shut Down T&S Rabbits" succeeded in closing down a network of rabbit meat and fur farms across the East Midlands region.
The specific future direction of cuniculture is unclear, but does not appear to be in danger of disappearing in any particular part of the world. The variety of applications, as well as the versatile utility of the species, appears sufficient to keep rabbit raising a going concern in one aspect or another around the planet.
Aspects of rabbit production
Meat rabbits
Rabbits have been raised for meat production in a variety of settings around the world. Smallholder or backyard operations remain common in many countries, while large-scale commercial operations are centered in Europe and Asia. For the smaller enterprise, multiple local rabbit breeds may be easier to use. Many local, "rustico", landrace or other heritage type breeds may be used only in a specific geographic area. Sub-par or "cull" animals from other breeding goals (laboratory, exhibition, show, wool, or pet) may also be used for meat, particularly in smallholder operations.
Counterintuitively, the giant rabbit breeds are rarely used for meat production, due to their extended growth rates (which lead to high feed costs) and their large bone size (which reduces the percentage of their weight that is usable meat). Dwarf breeds, too, are rarely used, due to the high production costs, slow growth, and low offspring rate.
In contrast to the multitude of breeds and types used in smaller operations, breeds such as the New Zealand and the Californian, along with hybrids of these breeds, are most frequently utilized for meat in commercial rabbitries. The primary qualities of good meat-rabbit breeding stock are growth rate and size at slaughter, but also good mothering ability. Specific lines of commercial breeds have been developed that maximize these qualities – rabbits may be slaughtered as early as seven weeks and does of these strains routinely raise litters of 8 to 12 kits. Other breeds of rabbit developed for commercial meat production include the Florida White and the Altex.
Rabbit breeding stock raised in France is particularly popular with meat rabbit farmers internationally, some being purchased as far away as China in order to improve the local rabbit herd.
Larger-scale operations attempt to maximize income by balancing land use, labor required, animal health, and investment in infrastructure. Specific infrastructure and strain qualities depend on the geographic area. An operation in an urban area may emphasize odor control and space utilization by stacking cages over each other with automatic cleaning systems that flush away faeces and urine. In rural sub-tropical and tropical areas, temperature control becomes more of an issue, and the use of air-conditioned buildings is common in many areas.
Breeding schedules for rabbits vary by individual operation. Prior to the development of modern balanced rabbit rations, rabbit breeding was limited by the nutrition available to the doe. Without adequate calories and protein, the doe would either not be fertile, would abort or resorb the foetuses during pregnancy, or would deliver small numbers of weak kits. Under these conditions, a doe would be re-bred only after weaning her last litter when the kits reached the age of two months. This allowed for a maximum of four litters per year. Advances in nutrition, such as those published by the USDA Rabbit Research Station, resulted in greater health for breeding animals and the survival of young stock. Likewise, offering superior, balanced nutrition to growing kits allowed for better health and less illness among slaughter animals. Current practices include the option of re-breeding the doe within a few days of delivery (closely matching the behavior of wild rabbits during the spring and early summer, when forage availability is at its peak.) This can result in up to eight or more litters annually. A doe of ideal meat-stock genetics can produce five times her body weight in fryers a year. Criticism of the more intensive breeding schedules has been made on the grounds that re-breeding that closely is excessively stressful for the doe. Determination of health effects of breeding schedules is made more difficult by the domestic rabbit's reproductive physiology – in contrast to several other mammal species, rabbits are more likely to develop uterine cancer when not used for breeding than when bred frequently.
In efficient production systems, rabbits can turn 20 percent of the proteins they eat into edible meat, compared to 22 to 23 percent for broiler chickens, 16 to 18 percent for pigs and 8 to 12 percent for beef; rabbit meat is more economical in terms of feed energy than beef.
"Rabbit fryers" are rabbits that are between 70 and 90 days old, weighing in live weight. "Rabbit roasters" are rabbits from 90 days to 6 months old, weighing in live weight. "Rabbit stewers" are rabbits 6 months or older, weighing over . "Dark fryers" (i.e., any color other than white) typically garner a lower price than "white fryers" (also called "albino fryers"), because of the slightly darker tinge to the meat. (Purely pink carcasses are preferred by most consumers.) Dark fryers are also harder to de-hide (skin) than white fryers.
In the United States, white fryers garner the highest prices per pound of live weight. In Europe, however, a sizable market remains for the dark fryers that come from older and larger rabbits. In the kitchen, dark fryers are typically prepared differently from white fryers.
In 1990, the world's annual production of rabbit meat was estimated to be 1.5 million tonnes. In 2014, the number was estimated at 2 million tonnes. China is among the world's largest producers and consumers of rabbit meat, accounting for some 30% of the world's total consumption. Within China itself, rabbits are raised in many provinces, with most of the rabbit meat (about 70% of the national production, equaling some 420,000 tonnes annually) being consumed in the Sichuan Basin (Sichuan Province and Chongqing), where it is particularly popular.
Well-known chef Mark Bittman wrote that domesticated rabbit "tastes like chicken", because both are "blank palettes on which we can layer whatever flavors we like".
Wool rabbits and pelt rabbits
Wool rabbits
Rabbits such as the Angora, American Fuzzy Lop, and Jersey Wooly produce wool. However, since the American Fuzzy Lop and Jersey Wooly are both dwarf breeds, only the much larger Angora breeds such as the English Angora, Satin Angora, Giant Angora, and French Angora are used for commercial wool production. Their long fur is sheared, combed, or plucked (gently pulling loose hairs from the body during molting) and then spun into yarn used to make a variety of products. Angora sweaters can be purchased in many clothing stores and is generally mixed with other types of wool. In 2010, 70% of Angora rabbit wool was produced in China. Rabbit wool, generically called Angora, is 5 times warmer than sheep's wool.
Pelt rabbits
A number of rabbit breeds have been developed with the fur trade in mind. Breeds such as the Rex, Satin, and Chinchilla are often raised for their fur. Each breed has fur characteristics and all have a wide range of colors and patterns. "Fur" rabbits are fed a diet especially balanced for fur production and the pelts are harvested when they have reached prime condition. Rabbit fur is widely used throughout the world. China imports much of its fur from Scandinavia (80%), and some from North America (5%), according to the USDA Foreign Agricultural Service GAIN Report CH7607.
Exhibition rabbits
Many rabbit keepers breed their rabbits for competition among other purebred rabbits of the same breed. Rabbits are judged according to the standards put forth by the governing associations of the particular country. These associations, being made up of people, may be distinctly political and reflect the preferences of particular persons on the governing boards. However, as mechanisms to preserve rare rabbit breeds, foster communication between breeders and encourage the education of the public, these organizations are invaluable. Examples include the American Rabbit Breeders Association and the British Rabbit Council.
Laboratory rabbits
Rabbits have been and continue to be used in laboratory work such as production of antibodies for vaccines and research of human male reproductive system toxicology. Experiments with rabbits date back to Louis Pasteur's work in France in the 1800s. In 1972, around 450,000 rabbits were used for experiments in the United States, decreasing to around 240,000 in 2006. The Environmental Health Perspective, published by the National Institute of Health, states, "The rabbit [is] an extremely valuable model for studying the effects of chemicals or other stimuli on the male reproductive system." According to the Humane Society of the United States, rabbits are also used extensively in the study of asthma, stroke prevention treatments, cystic fibrosis, diabetes, and cancer.
Rabbit cultivation intersects with research in two ways: first, the keeping and raising of animals for testing of scientific principles. Some experiments require the keeping of several generations of animals treated with a particular drug, in order to fully appreciate the side effects of that drug. There is also the matter of breeding and raising animals for experiments. The New Zealand White is one of the most commonly used breeds for research and testing. Specific strains of the New Zealand White have been developed, with differing resistance to disease and cancers. Additionally, some experiments call for the use of 'specific pathogen free' animals, which require specific husbandry and intensive hygiene.
Animal rights activists generally oppose animal experimentation for all purposes, and rabbits are no exception. The use of rabbits for the Draize test, which is used for, amongst other things, testing cosmetics on animals, has been cited as an example of cruelty in animal research. Albino rabbits are typically used in the Draize tests because they have less tear flow than other animals and the lack of eye pigment make the effects easier to visualize. Rabbits in captivity are uniquely subject to rabbitpox, a condition that has not been observed in the wild.
Husbandry
Modern methods for housing domestic rabbits vary from region to region around the globe and by type of rabbit, technological or financial opportunities and constraints, intended use, number of animals kept, and the particular preferences of the owner or farmer. Various goals include maximizing number of animals per land unit (especially common in areas with high land values or small living areas) minimizing labor, reducing cost, increasing survival and health of animals, and meeting specific market requirements (such as for clean wool, or rabbits raised on pasture). Not all of these goals are complementary. Where the keeping of rabbits has been regulated by governments, specific requirements have been put in place. Various industries also have commonly accepted practices which produce predictable results for that type of rabbit product.
Extensive cuniculture practices
Extensive cuniculture refers to the practice of keeping rabbits at a lower density and a lower production level than intensive culture. Specifically as relates to rabbits, this type of production was nearly universal prior to germ theory understanding of infectious parasites (especially coccidia) and the role of nutrition in prevention of abortion and reproductive loss. The most extensive rabbit "keeping" methods would be the harvest of wild or feral rabbits for meat or fur market, such as occurred in Australia prior to the 1990s. Warren-based cuniculture is somewhat more controlled, as the animals are generally kept to a specific area and a limited amount of supplemental feeding provided. Finally, various methods of raising rabbits with pasture as the primary food source have been developed. Pasturing rabbits within a fence (but not a cage), also known as colony husbandry, has not been commonly pursued due to the high death rate from weather and predators. More commonly (but still rare in terms of absolute numbers of rabbits and practitioners) is the practice of confining the rabbits to a moveable cage with an open or slatted floor so that the rabbits can access grass but still be kept at hand and protected from weather and predators. This method of growing rabbits does not typically result in an overall reduction for the need for supplemented feed. The growing period to market weight is much longer for grass fed rather than pellet fed animals, and many producers continue to offer small amounts of complete rations over the course of the growing period. Hutches or cages for this type of husbandry are generally made of a combination of wood and metal wire, made portable enough for a person to move the rabbits daily to fresh ground, and of a size to hold a litter of 6 to 12 rabbits at the market weight of . Protection from sun and driving rain are important health concerns, as is durability against predator attacks and the ability to be cleaned to prevent loss from coccidiosis.
Intensive cuniculture practices
Intensive cuniculture is more focused, efficient, and time-sensitive, utilizing a greater density of animals and higher turnover. The labor required to produce each harvested hide, kilogram of wool, or market fryer—and the quality thereof—may be higher or lower than for extensive methods. Successful operations raising healthy rabbits that produce durable goods range from thousands of animals to less than a dozen. Simple hutches, kitchen floors, and even natural pits may provide shelter from the elements, while the rabbits are fed from the garden or given what can be gathered as they grow to produce a community's foodstuffs and textiles. Intensive cuniculture can also be practiced in an enclosed, climate controlled barn where rows of cages house robust rabbits eating pellets and treats before a daily health inspection or weekly weight check. Veterinary specialists and biosecurity may be part of large-scale operations, while individual rabbits in smaller setups may receive better—or worse—care.
Challenges to successful production
Specific challenges to the keeping of rabbits vary by specific practices. Losses from coccidiosis are much more common when rabbits are kept on the ground (such as in warrens or colonies) or on solid floors than when in wire or slat cages that keep rabbits elevated away from urine and faeces. Pastured rabbits are more subject to predator attack. Rabbits kept indoors at an appropriate temperature rarely suffer heat loss compared to rabbits housed outdoors in summer. At the same time, if rabbits are housed inside without adequate ventilation, respiratory disease can be a significant cause of illness and death. Production does on fodder are rarely able to raise more than 3 litters a year without heavy losses from deaths of weak kits, abortion, and fetal resorption, all related to poor nutrition and inadequate protein intake. In contrast, rabbits fed commercial pelleted diets can face losses related to low fiber intake.
Exhibition and fancier societies
In the early 1900s, as animal fancy in general began to emerge, rabbit fanciers began to sponsor rabbit exhibitions and fairs in Western Europe and the United States. What became known as the "Belgian Hare Boom" began with the importation of the first Belgian Hares from England in 1888 and soon after the founding of the first rabbit club in America, the American Belgian Hare Association. From 1898 to 1901, many thousands of Belgian Hares were imported to America. Today, the Belgian Hare is considered one of the rarest breeds, with less than 200 in the United States as reported in a recent survey.
The American Rabbit Breeders Association (ARBA) was founded in 1910 and is the national authority on rabbit raising and rabbit breeds, having a uniform "Standard of Perfection", registration and judging system.
Conformation shows
Showing rabbits is an increasingly popular activity. Showing rabbits helps to improve the vigor and physical behavior of each breed through competitive selection. County fairs are common venues through which rabbits are shown in the United States. Rabbit clubs at local state and national levels hold many shows each year. Although only purebred animals are shown, a pedigree is not required to enter a rabbit in an ARBA-sanctioned show but is required to register the rabbit with ARBA. A rabbit must be registered in order to receive a Grand Champion certificate. Children's clubs such as 4H also include rabbit shows, usually in conjunction with county fairs. The ARBA holds an annual national convention which has as many as 25,000 animals competing from all over the world. The mega show moves to a different city each year. The ARBA also sponsors youth programs for families as well as underprivileged rural and inner city children to learn responsible care and breeding of domestic rabbits.
Genetics
The study of rabbit genetics is of interest to medical researchers, fanciers, and the fur and meat industries. Each of these groups has different needs for genetic information. In the biomedical research community and the pharmaceutical industry, rabbits genetics are important for producing antibodies, testing toxicity of consumer products, and in model organism research. Among rabbit fanciers and in the fiber and fur industry, the genetics of coat color and hair properties are paramount. The meat industry relies on genetics for disease resistance, feed conversion ratio, and reproduction potential.
The rabbit genome has been sequenced and is publicly available. The mitochondrial DNA has also been sequenced. In 2011, parts of the rabbit genome were re-sequenced in greater depth in order to expose variation within the genome.
Gene linkage maps
The early genetic research focused on linkage distance between various gross phenotypes using linkage analysis. Between 1924 and 1941, the relationship between c, y, b, du, En, l, r1, r2, A, dw, w, f, and br was established (phenotype is listed below).
c: albino
y: yellow fat
du: Dutch coloring
En: English coloring
l: angora
r1, r2: rex genes
A: Agouti
dw: dwarf gene
w: wide intermediate-color band
f: furless
br: brachydactyly
The distance between these genes is as follows, numbered by chromosome. The format is gene1—distance—gene2.
c — 14.4 — y — 28.4 — b
du — 1.2 — En — 13.1 — l
r1 — 17.2 — r2
A — 14.7 — dw — 15.4 — w
f — 28.3 — br
Color genes
There are 11 color gene groups (or loci) in rabbits. They are A, B, C, D, E, En, Du, P, Si, V, and W. Each locus has dominant and recessive genes. In addition to the loci there are also modifiers, which modify a certain gene. These include the rufous modifiers, color intensifiers, and plus/minus (blanket/spot) modifiers. A rabbit's coat has either two pigments (pheomelanin for yellow, and eumelanin for dark brown) or no pigment (for an albino rabbit).
Within each group, the genes are listed in order of dominance, with the most dominant gene first. In parentheses after the description is at least one example of a color that displays this gene.
Note: lower case are recessive and capital letters are dominant
"A" represents the agouti locus (multiple bands of color on the hair shaft). The genes are:
A: agouti ("wild color" or chestnut agouti, opal, chinchilla, etc.)
a(t): tan pattern (otter, tan, silver marten)
a: self- or non-agouti (black, chocolate)
"B" represents the brown locus. The genes are:
B: black (chestnut agouti, black otter, black)
b: brown (chocolate agouti, chocolate otter, chocolate)
"C" represents the color locus. The genes are:
C: full color (black)
c(ch3): dark chinchilla, removes yellow pigmentation (chinchilla, silver marten)
c(ch2): medium (light) chinchilla, slight reduction in eumelanin creating a more sepia tone in the fur rather than black.
c(ch1): light (pale) chinchilla (sable, sable point, smoke pearl, seal)
c(h): color sensitive expression of color. Warmer parts of the body do not express color. Known as Himalayan, the body is white with extremities (points) colored in black, blue, chocolate or lilac. Pink eyes.
c: albino (ruby-eyed white or REW)
"D" represents the dilution locus. This gene dilutes black to blue and chocolate to lilac.
D: dense color (chestnut agouti, black, chocolate)
d: diluted color (opal, blue or lilac)
"E" represents the extension locus. It works with the 'A' and 'C' loci and rufous modifiers. When it is recessive, it removes most black pigment. The genes are:
E(d): dominant black
E(s): steel (black removed from tips of fur, which then appear golden or silver)
E: normal
e(j): Japanese brindling (harlequin), black and yellow pigment broken into patches over the body. In a broken color pattern, this results in Tricolor.
e: most black pigment removed (agouti becomes red or orange, self- becomes tortoise)
"En" represents the plus/minus (blanket/spot) color locus. It is incompletely dominant and results in three possible color patterns:
EnEn: "Charlie" or a lightly marked broken with color on ears, on nose, and sparsely on body
Enen: "Broken" with roughly even distribution of color and white
enen: Solid color with no white areas
"Du" represents the Dutch color pattern (the front of the face, the front part of the body, and rear paws are white; the rest of the rabbit has colored fur). The genes are:
Du: absence of Dutch pattern
du(d): Dutch (dark)
du(w): Dutch (white)
"V" represents the vienna white locus. The genes are:
V: normal color
Vv: Vienna carrier; carries blue-eyed white gene. May appear as a solid color, with snips of white on nose and/or front paws, or Dutch marked.
v: vienna white (blue-eyed white or BEW)
"Si" represents the silver locus. The genes are:
Si: normal color
si: silver color (silver, silver fox)
"W" represents the middle yellow-white band locus and works with the agouti gene. The genes are:
W: normal width of yellow band
w: doubles yellow bandwidth (otter becomes tan, intensified red factors in Thrianta and Belgian Hare)
"P" represents the OCA type II form of albinism. P is used because it is an integral P protein mutation. The genes are:
P: normal color
p: albinism mutation. Removes eumelanin and causes pink eyes. (Will change, for example, a chestnut agouti into a shadow)
| Technology | Agriculture_2 | null |
170718 | https://en.wikipedia.org/wiki/Spruce | Spruce | A spruce is a tree of the genus Picea ( ), a genus of about 40 species of coniferous evergreen trees in the family Pinaceae, found in the northern temperate and boreal (taiga) regions of the Northern hemisphere. Picea is the sole genus in the subfamily Piceoideae. Spruces are large trees, from about 20 to 60 m (about 60–200 ft) tall when mature, and have whorled branches and conical form.
Spruces can be distinguished from other genera of the family Pinaceae by their needles (leaves), which are four-sided and attached singly to small persistent peg-like structures (pulvini or sterigmata) on the branches, and by their cones (without any protruding bracts), which hang downwards after they are pollinated. The needles are shed when 4–10 years old, leaving the branches rough with the retained pegs. In other similar genera, the branches are fairly smooth.
Spruce are used as food plants by the larvae of some Lepidoptera (moth and butterfly) species, such as the eastern spruce budworm. They are also used by the larvae of gall adelgids (Adelges species).
In the mountains of Dalarna, Sweden, scientists have found a Norway spruce, nicknamed Old Tjikko, which by reproducing through layering, has reached an age of 9,550 years and is claimed to be the world's oldest known living tree.
Description
Morphology
Determining that a tree is a spruce is not difficult; evergreen needles that are more or less quadrangled, and especially the pulvinus, give it away. Beyond that, determination can become more difficult. Intensive sampling in the Smithers/Hazelton/Houston area of British Columbia showed Douglas (1975), according to Coates et al. (1994), that cone scale morphology was the feature most useful in differentiating species of spruce; the length, width, length: width ratio, the length of free scale (the distance from the imprint of the seed wing to the tip of the scale), and the percentage free scale (length of free scale as a percentage of the total length of the scale) were most useful in this regard. Daubenmire (1974), after range-wide sampling, had already recognized the importance of the two latter characters. Taylor (1959) had noted that the most obvious morphological difference between typical Picea glauca and typical P. engelmannii was the cone scale, and Horton (1956,1959) found that the most useful diagnostic features of the two spruces are in the cone; differences occur in the flower, shoot and needle, "but those in the cone are most easily assessed" (Horton 1959). Coupé et al. (1982) recommended that cone scale characters be based on samples taken from the midsection of each of ten cones from each of five trees in the population of interest.
Without cones, morphological differentiation among spruce species and their hybrids is more difficult. Species classification for seeds collected from spruce stands in which introgressive hybridization between white and Sitka spruces (P. sitchensis) may have occurred is important for determining appropriate cultural regimens in the nursery. If, for instance, white spruce grown at container nurseries in southwestern British Columbia are not given an extended photoperiod, leader growth ceases early in the first growing season, and seedlings do not reach the minimum height specifications. But, if an extended photoperiod is provided for Sitka spruce, seedlings become unacceptably tall by the end of the first growing season. Species classification of seedlots collected in areas where hybridization of white and Sitka spruces has been reported has depended on (i) easily measured cone scale characters of seed trees, especially free scale length, (ii) visual judgements of morphological characters, e.g., growth rhythm, shoot and root weight, and needle serration, or (iii) some combination of (i) and (ii) (Yeh and Arnott 1986). Useful to a degree, these classification procedures have important limitations; genetic composition of the seeds produced by a stand is determined by both the seed trees and the pollen parents, and species classification of hybrid seedlots and estimates of their level of introgression on the basis of seed-tree characteristics can be unreliable when hybrid seedlots vary in their introgressiveness in consequence of spatial and temporal variations in contributions from the pollen parent (Yeh and Arnott 1986). Secondly, morphological characters are markedly influenced by ontogenetic and environmental influences, so that to discern spruce hybrid seedlot composition with accuracy, hybrid seedlots must differ substantially in morphology from both parent species. Yeh and Arnott (1986) pointed out the difficulties of estimating accurately the degree of introgression between white and Sitka spruces; introgression may have occurred at low levels, and/or hybrid seed lots may vary in their degree of introgression in consequence of repeated backcrossing with parental species.
Growth
Spruce seedlings are most susceptible immediately following germination, and remain highly susceptible through to the following spring. More than half of spruce seedling mortality probably occurs during the first growing season and is also very high during the first winter, when seedlings are subjected to freezing damage, frost heaving and erosion, as well as smothering by litter and snow-pressed vegetation. Seedlings that germinate late in the growing season are particularly vulnerable because they are tiny and have not had time to harden off fully.
Mortality rates generally decrease sharply thereafter, but losses often remain high for some years. "Establishment" is a subjective concept based on the idea that once a seedling has successfully reached a certain size, not much is likely to prevent its further development. Criteria vary, of course, but Noble and Ronco (1978), for instance, considered that seedlings four to five years old, or 8 cm to 10 cm tall, warranted the designation "established", since only unusual factors such as snow mold, fire, trampling, or predation would then impair regeneration success. Eis (1967) suggested that in dry habitats on either mineral soil or litter seedbeds a 3-year-old seedling may be considered established; in moist habitats, seedlings may need 4 or 5 years to become established on mineral soil, possibly longer on litter seedbeds.
Growth remains very slow for several to many years. Three years after shelterwood felling in subalpine Alberta, dominant regeneration averaged 5.5 cm in height in scarified blocks, and 7.3 cm in non-scarified blocks (Day 1970), possibly reflecting diminished fertility with the removal of the A horizon.
Taxonomy
Classification
DNA analyses have shown that traditional classifications based on the morphology of needle and cone are artificial. A 2006 study found that P. breweriana had a basal position, followed by P. sitchensis, and the other species were further divided into three clades, suggesting that Picea originated in North America. The oldest record of spruce that has been found in the fossil record is from the Early Cretaceous (Valanginian) of western Canada, around 136 million years old.
Species
, Plants of the World Online accepted 37 species. The grouping is based on Ran et al. (2006).
Basal species:
Picea breweriana – Brewer's spruce, Klamath Mountains, North America; local endemic
Picea sitchensis – Sitka spruce, Pacific coast of North America; the largest species, to 95 m tall; important in forestry
Clade I (Northern and western North America, in boreal forests or high mountains)
Picea engelmannii – Engelmann spruce, western North American mountains; important in forestry
Picea glauca, syn. Picea laxa – white spruce, northern North America; important in forestry
Clade II (throughout Asia, mostly in mountainous areas, a few isolated populations in higher elevations of Mexico,)
Picea brachytyla – Sargent's spruce, southwest China
Picea chihuahuana – Chihuahua spruce, northwest Mexico (rare)
Picea farreri – Burmese spruce, northeast Burma, southwest China (mountains)
Picea likiangensis – Likiang spruce, southwest China
Picea martinezii – Martinez spruce, northeast Mexico (very rare, endangered)
Picea maximowiczii – Maximowicz spruce, Japan (rare, mountains)
Picea morrisonicola – Taiwan spruce, Taiwan (high mountains)
Picea neoveitchii – Veitch's spruce, northwest China (rare, endangered)
Picea orientalis – Caucasian spruce or Oriental spruce, Caucasus, northeast Turkey
Picea polita, syn. Picea torano – tiger-tail spruce, Japan
Picea purpurea – purple cone spruce, western China
Picea schrenkiana – Schrenk's spruce, mountains of central Asia
Picea smithiana – morinda spruce, western Himalaya, eastern Afghanistan, northern and northwest India
Picea spinulosa – Sikkim spruce, northeast India (Sikkim), eastern Himalaya
Picea wilsonii – Wilson's spruce, western China
Clade III (Europe, Asia, and North America, mostly in boreal forests or mountainous areas)
Picea abies – Norway spruce, Europe; important in forestry, the original Christmas tree
Picea alcoquiana – ("P. bicolor") Alcock's spruce, central Japan (mountains)
Picea asperata – dragon spruce, western China; several varieties
Picea crassifolia – Qinghai spruce, China
Picea glehnii – Glehn's spruce, northern Japan, Sakhalin
Picea jezoensis – Jezo spruce, northeast Asia, Kamchatka south to Japan
Picea koraiensis – Korean spruce, Korea, northeast China
Picea koyamae – Koyama's spruce, Japan (mountains)
Picea mariana – black spruce, northern North America
Picea meyeri – Meyer's spruce, northern China (from Inner Mongolia to Gansu)
Picea obovata – Siberian spruce, north Scandinavia, Siberia; often treated as a variant of P. abies (and hybridises with it), but has distinct cones
Picea omorika – Serbian spruce, Serbia and Bosnia; local endemic; important in horticulture
Picea pungens – blue spruce or Colorado spruce, Rocky Mountains, North America; important in horticulture
Picea retroflexa – green dragon spruce, China
Picea rubens – red spruce, northeastern North America; important in forestry, known as Adirondack in musical-instrument making
Others
Picea aurantiaca Mast.
Picea austropanlanica Silba
Picea linzhiensis (W.C.Cheng & L.K.Fu) Rushforth
Hybrids
Picea × albertiana S.Br.
Picea × fennica (Regel) Kom.
Picea × lutzii Little
[[Picea × notha|Picea × notha]] Rehder
Fossil species
There are also a number of extinct species identified from fossil evidence:
Cones
†Picea anadyrensis
†Picea antiqua
†Picea beckii
†Picea bilibinii
†Picea burtonii
†Picea camtschatica
†Picea deweyensis
†Picea diettertiana
†Picea eichhornii
†Picea evenica
†Picea fimbriata
†Picea garoensis
†Picea harrimani
†Picea hondoensis
†Picea indigirensis
†Picea latibracteata
†Picea latisquamosa
†Picea metechensis
†Picea mioorientalis
†Picea mugodzharica
†Picea oligocaenica
†Picea pevekensis
†Picea praeajanensis
†Picea protopicea
†Picea sookensis
†Picea sugaii
†Picea suifunensis
†Picea vassiljevii
†Picea vitjasii
†Picea wolfei
†Picea wollosowiczii
†Picea yanensis
Foliage
†Picea cretacea
†Picea echinata
†Picea korfiensis
†Picea lahontensis
†Picea morosovae
†Picea nakauchii
†Picea palaeomorika
†Picea quilchensis
†Picea tranquillensis
Pollen
†Picea alata
†Picea bella
†Picea complanatiformis
†Picea depressa
†Picea distorta
†Picea distracta
†Picea exilioides
†Picea gigantissima
†Picea grandipollinia
†Picea grandis
†Picea grandivescipites
†Picea kryshtofovichii
†Picea longisaccata
†Picea media
†Picea mesophytica
†Picea minor
†Picea multigruma
†Picea omoriciformis
†Picea parvireticulata
†Picea pseudorotundiformis
†Picea rara
†Picea sacculifera
†Picea samoilovitchiana
†Picea schrenkianiformis
†Picea scotica
†Picea singularis
†Picea spirelliformis
†Picea sutschanensis
†Picea tasaranica
†Picea tobolica
†Picea valanjinica
Seeds
†Picea altaica
†Picea hiyamensis
†Picea kaneharae
†Picea kanoi
†Picea magna
†Picea pinifructus
†Picea sonomensis
Wood
†Picea palaeomaximowiczii
†Picea wakimizui
†Picea withamii
Multiple organs
†Picea columbiensis
†Picea critchfieldii - Late Quaternary North America.
†Picea farjonii
†Picea heisseana
†Picea koribae
†Picea miocenica
†Picea rotundosquamosa
†Picea ugoana
Unspecified
†Picea alba
†Picea albertensis
†Picea snatolensis
Fossil species formerly placed in Picea
†Picea cliffwoodensis moved to †Pityostrobus cliffwoodensis
†Picea succinifera moved to Pinus succinifera
Cultivation
In the realm of spruce trees, the presence of Dendroctonus micans beetles significantly impacts their health and vitality. These beetles, particularly the males, display territorial behavior, diligently defending areas that are attractive to females for mating. By safeguarding these regions and providing suitable host trees, they create an environment conducive to egg-laying, thereby ensuring their reproductive success. This territoriality is closely tied to male reproductive prowess and plays a crucial role in understanding the dynamics of beetle populations and their impact on the overall health of spruce trees. Moreover, the home range of Dendroctonus micans varies based on the availability and density of host trees, with individual beetles dispersing across their habitat in search of optimal nesting sites.
Etymology Spruce, (1412), and (1378) seem to have been generic terms for commodities brought to England by Hanseatic merchants (especially beer, boards, wooden chests and leather), and the tree thus was believed to be particular to Prussia, which for a time was figurative in England as a land of luxuries. It can be argued that the word is actually derived from the Old French term , meaning literally Prussia.
Ecology
Diseases
Sirococcus blight (Deuteromycotina, Coelomtcetes)
The closely related species Sirococcus conigenus and Sirococcus piceicola cause shoot blight and seedling mortality of conifers in North America, Europe, and North Africa. Twig blight damage to seedlings of white and red spruces in a nursery near Asheville, North Carolina, was reported by Graves (1914). Hosts include white, black, Engelmann, Norway, and red spruces, although they are not the plants most commonly damaged. Sirococcus blight of spruces in nurseries show up randomly in seedlings to which the fungus was transmitted in infested seed. First-year seedlings are often killed, and larger plants may become too deformed for planting. Outbreaks involving < 30% of spruce seedlings in seedbeds have been traced to seed lots in which only 0.1% to 3% of seeds were infested. Seed infestation has in turn been traced to the colonization of spruce cones by S. conigenus in forests of the western interior. Infection develops readily if conidia are deposited on succulent plant parts that remain wet for at least 24 hours at 10 °C to 25 °C. Longer periods of wetness favour increasingly severe disease. Twig tips killed during growth the previous year show a characteristic crook.
Rhizosphaera kalkhoffi needle cast Rhizosphaera infects white spruce, blue spruce (Picea pungens), and Norway spruces throughout Ontario, causing severe defoliation and sometimes killing small, stressed trees. White spruce is intermediately susceptible. Dead needles show rows of black fruiting bodies. Infection usually begins on lower branches. On white spruce, infected needles are usually retained on the tree into the following summer. The fungicide Chlorthalonil is registered for controlling this needle cast (Davis 1997).
Valsa kunzei branch and stem canker
A branch and stem canker associated with the fungus Valsa kunzei Fr. var. picea was reported on white and Norway spruces in Ontario (Jorgensen and Cafley 1961) and Quebec (Ouellette and Bard 1962). In Ontario, only trees of low vigour were affected, but in Quebec vigorous trees were also infected.
Predators
Small mammals ingest conifer seeds, and also consume seedlings. Cage feeding of deer mice (Peromyscus maniculatus) and red-backed vole (Myodes gapperi) showed a daily maximum seed consumption of 2000 white spruce seeds and of 1000 seeds of lodgepole pine, with the two species of mice consuming equal amounts of seed, but showing a preference for the pine over the spruce (Wagg 1963). The short-tailed meadow vole (Microtus pennsylvanicus Ord) voraciously ate all available white spruce and lodgepole pine seedlings, pulling them out of the ground and holding them between their front feet until the whole seedling had been consumed. Wagg (1963) attributed damage observed to the bark and cambium at ground level of small white spruce seedlings over several seasons to meadow voles.
Once shed, seeds contribute to the diet of small mammals, e.g., deer mice, red-backed voles, mountain voles (Microtus montanus), and chipmunks (Eutamias minimus). The magnitude of the loss is difficult to determine, and studies with and without seed protection have yielded conflicting results. In western Montana, for example, spruce seedling success was little better on protected than on unprotected seed spots (Schopmeyer and Helmers 1947), but in British Columbia spruce regeneration depended on protection from rodents (Smith 1955).
An important albeit indirect biotic constraint on spruce establishment is the depredation of seed by squirrels. As much as 90% of a cone crop has been harvested by red squirrels (Zasada et al. 1978). Deer mice, voles, chipmunks, and shrews can consume large quantities of seed; one mouse can eat 2000 seeds per night. Repeated applications of half a million seeds/ha failed to produce the 750 trees/ha sought by Northwest Pulp and Power, Ltd., near Hinton, Alberta (Radvanyi 1972), but no doubt left a lot of well-fed small mammals. Foraging by squirrels for winter buds (Rowe 1952) has not been reported in relation to young plantations, but Wagg (1963) noted that at Hinton AB, red squirrels were observed cutting the lateral and terminal twigs and feeding on the vegetative and flower buds of white spruce.
Red squirrels in Alaska have harvested as much as 90% of a cone crop (Zasada et al. 1978); their modus operandi is to cut off great numbers of cones with great expedition early in the fall, and then "spend the rest of the fall shelling out the seeds". In Manitoba, Rowe (1952) ascribed widespread severing of branch tips 5 cm to 10 cm long on white spruce ranging "from sapling to veteran size" to squirrels foraging for winter buds, cone failure having excluded the more usual food source. The damage has not been reported in relation to small trees, outplants or otherwise.
Porcupines (Erethizon dorsatum L.) may damage spruce (Nienstaedt 1957), but prefer red pine. Bark-stripping of white spruce by black bear (Euarctos americanus perniger) is locally important in Alaska (Lutz 1951), but the bark of white spruce is not attacked by field mice (Microtus pennsylvanicus Ord), even in years of heavy infestation.
Pests
The eastern spruce budworm (Choristoneura fumiferana) is a major pest of spruce trees in forests throughout Canada and the eastern United States. Two of the main host plants are black spruce and white spruce. Population levels oscillate, sometimes reaching extreme outbreak levels that can cause extreme defoliation of and damage to spruce trees. To reduce destruction, there are multiple methods of control in place, including pesticides.
Horntails, or Wood Wasps, use this tree for egg laying and the larvae will live in the outer inch of the tree under the bark.
Spruce beetles (Dendroctonus rufipennis) have destroyed swathes of spruce forest in western North America from Alaska to Wyoming.
Uses
Timber
Spruce is useful as a building wood, commonly referred to by several different names including North American timber, SPF (spruce, pine, fir) and whitewood (the collective name for spruce wood). It is commonly used in Canadian Lumber Standard graded wood. Spruce wood is used for many purposes, ranging from general construction work and crates to highly specialised uses in wooden aircraft. The Wright brothers' first aircraft, the Flyer, was built of spruce.
Because this species has no insect or decay resistance qualities after logging, it is generally recommended for construction purposes as indoor use only (indoor drywall framing, for example). Spruce wood, when left outside cannot be expected to last more than 12–18 months depending on the type of climate it is exposed to.
Pulpwood
Spruce is one of the most important woods for paper uses, as it has long wood fibres which bind together to make strong paper. The fibres are thin walled and collapse to thin bands upon drying. Spruces are commonly used in mechanical pulping as they are easily bleached. Together with northern pines, northern spruces are commonly used to make NBSK. Spruces are cultivated over vast areas as pulpwood.
Food and medicine
The fresh shoots of many spruces are a natural source of vitamin C. Captain Cook made alcoholic sugar-based spruce beer during his sea voyages in order to prevent scurvy in his crew. The leaves and branches, or the essential oils, can be used to brew spruce beer.
In Finland, young spruce buds are sometimes used as a spice, or boiled with sugar to create spruce bud syrup. In survival situations spruce needles can be directly ingested or boiled into a tea. This replaces large amounts of vitamin C. Also, water is stored in a spruce's needles, providing an alternative means of hydration . Spruce can be used as a preventive measure for scurvy in an environment where meat is the only prominent food source .
Tonewood
Spruce is the standard material used in soundboards for many musical instruments, including guitars, mandolins, cellos, violins, and the soundboard at the heart of a piano and the harp.
Wood used for this purpose is referred to as tonewood.
Spruce, along with cedar, is often used for the soundboard/top of an acoustic guitar. The main types of spruce used for this purpose are Sitka, Engelmann, Adirondack and European spruces.
Other uses
The resin was used in the manufacture of pitch in the past (before the use of petrochemicals); the scientific name Picea derives from Latin "pitch pine" (referring to Scots pine), from , an adjective from "pitch".
Native Americans in North America use the thin, pliable roots of some species for weaving baskets and for sewing together pieces of birch bark for canoes. | Biology and health sciences | Gymnosperms | null |
170722 | https://en.wikipedia.org/wiki/Fir | Fir | Firs are evergreen coniferous trees belonging to the genus Abies () in the family Pinaceae. There are approximately 48–65 extant species, found on mountains throughout much of North and Central America, Eurasia, and North Africa. The genus is most closely related to Keteleeria, a small genus confined to eastern Asia.
The genus name is derived from the Latin "to rise" in reference to the height of its species. The common English name originates with the Old Norse fyri or the Old Danish fyr.
They are large trees, reaching heights of tall with trunk diameters of when mature. Firs can be distinguished from other members of the pine family by the way in which their needle-like leaves are attached singly to the branches with a base resembling a suction cup, and by their cones, which, like those of cedars, stand upright on the branches like candles and disintegrate at maturity.
Identification of the different species is based on the size and arrangement of the leaves, the size and shape of the cones, and whether the bract scales of the cones are long and exserted, or short and hidden inside the cone.
Description
Leaves
Firs can be distinguished from other members of the pine family by the unique attachment of their needle-like leaves to the twig by a base that resembles a small suction cup. The leaves are significantly flattened, sometimes even looking like they are pressed, as in A. sibirica.
The leaves have two whitish lines on the bottom, each of which is formed by wax-covered stomatal bands. In most species, the upper surface of the leaves is uniformly green and shiny, without stomata or with a few on the tip, visible as whitish spots. Other species have the upper surface of leaves dull, greyish green or bluish to silvery (glaucous), coated by wax with variable number of stomatal bands, and not always continuous. An example species with shiny green leaves is A. alba, and an example species with matt waxy leaves is A. concolor.
The tips of leaves are usually more or less notched (as in A. firma), but sometimes rounded or dull (as in A. concolor, A. magnifica) or sharp and prickly (as in A. bracteata, A. cephalonica, A. holophylla). The leaves of young plants are usually sharper.
The leaves are arranged spirally on the shoots, but by being twisted at their base, the way they spread from the shoot is diverse; in some species comb-like ('pectinate'), with the leaves flat on either side of the shoot (e.g. A. alba, A. grandis), in others, the leaves remain radial (e.g. A. pinsapo)
Foliage in the upper crown on cone-bearing branches is different, with the leaves shorter, curved, and sometimes sharp.
Cones
Firs differ from other conifers in having erect, cylindrical cones long that disintegrate at maturity to release the winged seeds. In contrast to spruces, fir cones are erect; they do not hang, unless heavy enough to twist the branch with their weight.
The mature cones are usually brown. When young in summer, they can be green:
A. grandis, A. holophylla
or reddish:
A. alba, A. cephalonica, A. nordmanniana
or bloomed pale glaucous or pinkish:
A. numidica, A. pinsapo
or purple to blue, sometimes very dark blue, almost black:
A. forrestii, A. fraseri, A. homolepis, A. lasiocarpa, A. pindrow.
Many species are polymorphic in cone colour, with different individuals of the same species producing either green or purple cones:
A. concolor, A. koreana (usually purple, rarely green, such as the cultivar 'Flava'), A. magnifica (usually green, occasionally purple), A. nephrolepis (f. chlorocarpa green), A. sibirica, A. veitchii (f. olivacea green)
The cone scale bracts can be short and hidden in the mature cone, or long and exposed ('exserted'); this can vary even within a species, e.g. in Abies magnifica var. magnifica, the bracts are hidden, but in var. critchfieldii and var. shastensis, they are exserted. The bracts scales are often a different colour to the cone scales, which can make for a very attractive combination valued in ornamental trees.
Classification
The oldest pollen assignable to the genus dates to the Late Cretaceous in Siberia, with records of leaves and reproductive organs across the Northern Hemisphere from the Eocene onwards.
Section Abies
Section Abies is found in central, south, and eastern Europe and Asia Minor.
Abies alba – silver fir or European silver fir
Abies nebrodensis – Sicilian fir
Abies borisii-regis – Bulgarian fir
Abies cephalonica – Greek fir
Abies nordmanniana – Caucasian fir or Nordmann fir
Abies nordmanniana subsp. equi-trojani – Kazdağı fir, Turkish fir
Abies pinsapo – Spanish fir
Abies pinsapo var. marocana – Moroccan fir
Abies numidica – Algerian fir
Abies cilicica – Syrian fir
Section Balsamea
Section Balsamea is found in northern Asia and North America, and high mountains further south.
Abies fraseri – Fraser's fir
Abies balsamea – balsam fir
Abies balsamea var. phanerolepis – bracted balsam fir
Abies lasiocarpa – subalpine fir
Abies lasiocarpa var. arizonica – corkbark fir
Abies lasiocarpa var. bifolia – Rocky Mountains subalpine fir
Abies sibirica – Siberian fir
Abies sibirica var. semenovii
Abies sachalinensis – Sakhalin fir
Abies koreana – Korean fir
Abies nephrolepis – Khinghan fir
Abies veitchii – Veitch's fir
Abies veitchii var. sikokiana – Shikoku fir
Section Grandis
Section Grandis is found in western North America to Mexico, Guatemala, Honduras and El Salvador, in lowlands in the north, moderate altitudes in south.
Abies grandis – grand fir or giant fir
Abies grandis var. grandis – Coast grand fir
Abies grandis var. idahoensis – interior grand fir
Abies concolor – white fir
Abies concolor subsp. concolor – Rocky Mountain white fir or Colorado white fir
Abies concolor subsp. lowiana – Low's white fir or Sierra Nevada white fir
Abies durangensis – Durango fir
Abies durangensis var. coahuilensis – Coahuila fir
Abies flinckii – Jalisco fir
Abies guatemalensis – Guatemalan fir
Abies guatemalensis var. guatemalensis
Abies guatemalensis var. jaliscana
Abies vejarii
Section Momi
Section Momi is found in east and central Asia and the Himalaya, generally at low to moderate altitudes.
Abies kawakamii – Taiwan fir
Abies homolepis – Nikko fir
Abies recurvata – Min fir
Abies recurvata var. ernestii – Min fir
Abies firma – Momi fir
Abies beshanzuensis – Baishanzu fir
Abies holophylla – Manchurian fir
Abies chensiensis – Shensi fir
Abies chensiensis subsp. salouenensis – Salween fir
Abies pindrow – Pindrow fir
Abies ziyuanensis – Ziyuan fir
Section Amabilis
Section Amabilis is found in the Pacific Coast mountains in North America and Japan, in high rainfall areas.
Abies amabilis – Pacific silver fir
Abies mariesii – Maries' fir
Section Pseudopicea
Section Pseudopicea is found in the Sino – Himalayan mountains at high altitudes.
Abies delavayi – Delavay's fir
Abies delavayi var. nukiangensis
Abies delavayi var. motuoensis
Abies delavayi subsp. fansipanensis
Abies fabri – Faber's fir
Abies fabri subsp. minensis
Abies forrestii – Forrest's fir
Abies densa – Bhutan fir
Abies spectabilis – East Himalayan fir
Abies fargesii – Farges' fir
Abies fanjingshanensis – Fanjingshan fir
Abies yuanbaoshanensis – Yuanbaoshan fir
Abies squamata – flaky fir
Section Oiamel
Section Oiamel is found in central Mexico at high altitudes.
Abies religiosa – sacred fir
Abies hickelii – Hickel's fir
Abies hickelii var. oaxacana – Oaxaca fir
Section Nobilis
Section Nobilis (western U.S., high altitudes)
Abies procera – noble fir
Abies magnifica – red fir
Abies magnifica var. shastensis – Shasta red fir
Section Bracteata
Section Bracteata (California coast)
Abies bracteata – bristlecone fir
?†Abies rigida - (Priabonian-Chattian; Colorado)
Section Incertae sedis
Section Incertae sedis
†Abies milleri – (Extinct) Early Eocene
Ecology
Firs are used as food plants by the caterpillars of some Lepidoptera species, including Chionodes abella (recorded on white fir), autumnal moth, conifer swift (a pest of balsam fir), the engrailed, grey pug, mottled umber, pine beauty and the tortrix moths Cydia illutana (whose caterpillars are recorded to feed on European silver fir cone scales) and C. duplicana (on European silver fir bark around injuries or canker).
Abies religiosa (sacred fir) trees give roosting shelter to overwintering monarch butterflies.
Phytochemistry
Abies produce a variety of terpenoids. The analyses of the Zavarin groupfrom Smedman et al. 1969 to Zavarin et al. 1977showed variation in terpenoid composition of the bark by genetics, geography, age and size of the tree.
Uses
Wood of most firs is considered unsuitable for general timber use and is often used as pulp or for the manufacture of plywood and rough timber. It is commonly used in Canadian Lumber Standard graded wood. Because this genus has no insect or decay resistance qualities after logging, it is generally recommended in construction purposes for indoor use only (e.g. indoor drywall on framing). Firwood left outside cannot be expected to last more than 12 to 18 months, depending on the type of climate it is exposed to.
Caucasian fir, noble fir, Fraser's fir and balsam fir are popular Christmas trees, generally considered to be the best for this purpose, with aromatic foliage that does not shed many needles on drying out. Many are also decorative garden trees, notably Korean fir and Fraser's fir, which produce brightly coloured cones even when very young, still only tall. Many fir species are grown in botanic gardens and other specialist tree collections in Europe and North America.
Abies spectabilis or Talispatra is used in Ayurveda as an antitussive (cough suppressant) drug.
| Biology and health sciences | Gymnosperms | null |
170742 | https://en.wikipedia.org/wiki/Betulaceae | Betulaceae | Betulaceae, the birch family, includes six genera of deciduous nut-bearing trees and shrubs, including the birches, alders, hazels, hornbeams, hazel-hornbeam, and hop-hornbeams, numbering a total of 167 species. They are mostly natives of the temperate Northern Hemisphere, with a few species reaching the Southern Hemisphere in the Andes in South America. Their typical flowers are catkins and often appear before leaves.
In the past, the family was often divided into two families, Betulaceae (Alnus, Betula) and Corylaceae (the rest). Recent treatments, including the Angiosperm Phylogeny Group, have described these two groups as subfamilies within an expanded Betulaceae: Betuloideae and Coryloideae.
Betulaceae flowers are monoecious, meaning that they have both male and female flowers on the same tree. Their flowers present as catkins and are small and inconspicuous, often with reduced perianth parts. These flowers have large feathery stamen and produce a high volume of pollen, as they rely on wind pollination. Their leaves are simple, with alternate arrangement and doubly serrate margins.
Evolutionary history
The Betulaceae are believed to have originated at the end of the Cretaceous period (about 70 million years ago) in central China. This region at the time would have had a Mediterranean climate due to the proximity of the Tethys Sea, which covered parts of present-day Tibet and Xinjiang into the early Tertiary period. This point of origin is supported by the fact that all six genera and 52 species are native to this region, many of those being endemic. All six modern genera are believed to have diverged fully by the Oligocene, with all genera in the family (with the exception of Ostryopsis) having a fossil record stretching back at least 20 million years from the present.
According to molecular phylogeny, the closest relatives of the Betulaceae are the Casuarinaceae, or the she-oaks.
Uses
The common hazel (Corylus avellana) and the filbert (Corylus maxima) are important orchard plants, grown for their edible nuts.
The other genera include a number of popular ornamental trees, widely planted in parks and large gardens; several of the birches are particularly valued for their smooth, brightly coloured bark.
The wood is generally hard, tough and heavy, hornbeams particularly so; several species were of significant importance in the past where very hard wood capable of withstanding heavy wear was required, such as for cartwheels, water wheels, cog wheels, tool handles, chopping boards, and wooden pegs. In most of these uses, wood has now been replaced by metal or other man-made materials.
Subfamilies and genera
Extant genera
Betuloideae
Alnus Mill. 1754 – alder
Betula L. 1753 – birch
Coryloideae
Carpinus L. 1753 – hornbeam
Corylus L. 1753 – hazel
Ostrya Scop. 1760 – hop-hornbeam
Ostryopsis Decne. 1873 – hazel-hornbeam
Fossil genera
Betuloidea
Alnipollenites (pollen)
†Alnoxylon (wood)
†Alnuspollenites (pollen)
†Betulapollenites (pollen)
†Betulaepollenites (pollen)
†Paralnoxylon (wood)
Coryloidea
†Asterocarpinus (fruits)
†Coryloides (fruits)
†Cranea (fruits, flowers, pollen & Leaves)
†Corylocarpinus (fruits)
†Craspedodromophyllum (leaves)
†Kardiasperma (fruits)
†Palaeocarpinus (fruits)
†Paracarpinus (leaves)
Incertae sedis
†Alniphyllum (leaves)
†Alnites (leaves, flowers, fruits)
†Alnophyllum (leaves)
†Betulites (flowers)
†Betulacites (pollen)
†Betulinium (wood)
†Trivestibulopollenites (pollen)
†Betuliphyllum (leaves)
†Betuloidites (pollen)
†Betuloxylon (wood)
†Carpinicarpus (fruits)
†Carpiniphyllum (leaves)
†Carpinites (flowers)
†Carpinoxylon (wood)
†Carpinuspollenites (pollen)
†Castanopsispollenites (pollen)
†Clethrites (wood)
†Corylipollenites (pollen)
†Corylites (leaves)
†Coryloidites (pollen)
†Corylophyllum (leaves)
†Coryloxylon (wood)
†Coryluspollenites (pollen)
†Eucarpinoxylon (wood)
†Ostryoipollenites (pollen)
†Phegites (wood)
†Rhizoalnoxylon (root)
†Triporopollenites (pollen)
Phylogenetic systematics
Modern molecular phylogenetics suggest the following relationships:
| Biology and health sciences | Fagales | Plants |
170757 | https://en.wikipedia.org/wiki/Cyclotron%20radiation | Cyclotron radiation | In particle physics, cyclotron radiation is electromagnetic radiation emitted by non-relativistic accelerating charged particles deflected by a magnetic field. The Lorentz force on the particles acts perpendicular to both the magnetic field lines and the particles' motion through them, creating an acceleration of charged particles that causes them to emit radiation as a result of the acceleration they undergo as they spiral around the lines of the magnetic field.
The name of this radiation derives from the cyclotron, a type of particle accelerator used since the 1930s to create highly energetic particles for study. The cyclotron makes use of the circular orbits that charged particles exhibit in a uniform magnetic field. Furthermore, the period of the orbit is independent of the energy of the particles, allowing the cyclotron to operate at a set frequency. Cyclotron radiation is emitted by all charged particles travelling through magnetic fields, not just those in cyclotrons. Cyclotron radiation from plasma in the interstellar medium or around black holes and other astronomical phenomena is an important source of information about distant magnetic fields.
Properties
The power (energy per unit time) of the emission of each electron can be calculated:
where E is energy, t is time, is the Thomson cross section (total, not differential), B is the magnetic field strength, v is the velocity perpendicular to the magnetic field, c is the speed of light and is the permeability of free space.
Cyclotron radiation has a spectrum with its main spike at the same fundamental frequency as the particle's orbit, and harmonics at higher integral factors. Harmonics are the result of imperfections in the actual emission environment, which also create a broadening of the spectral lines. The most obvious source of line broadening is non-uniformities in the magnetic field; as an electron passes from one area of the field to another, its emission frequency will change with the strength of the field. Other sources of broadening include collisional broadening as the electron will invariably fail to follow a perfect orbit, distortions of the emission caused by interactions with the surrounding plasma, and relativistic effects if the charged particles are sufficiently energetic. When the electrons are moving at relativistic speeds, cyclotron radiation is known as synchrotron radiation.
The recoil experienced by a particle emitting cyclotron radiation is called radiation reaction. Radiation reaction acts as a resistance to motion in a cyclotron; and the work necessary to overcome it is the main energetic cost of accelerating a particle in a cyclotron. Cyclotrons are prime examples of systems which experience radiation reaction.
Examples
In the context of magnetic fusion energy, cyclotron radiation losses translate into a requirement for a minimum plasma energy density in relation to the magnetic field energy density.
Cyclotron radiation would likely be produced in a high altitude nuclear explosion. Gamma rays produced by the explosion would ionize atoms in the upper atmosphere and those free electrons would interact with the Earth's magnetic field to produce cyclotron radiation in the form of an electromagnetic pulse (EMP). This phenomenon is of concern to the military as the EMP may damage solid state electronic equipment.
| Physical sciences | Electromagnetic radiation | Physics |
170808 | https://en.wikipedia.org/wiki/Synchrotron%20radiation | Synchrotron radiation | Synchrotron radiation (also known as magnetobremsstrahlung) is the electromagnetic radiation emitted when relativistic charged particles are subject to an acceleration perpendicular to their velocity (). It is produced artificially in some types of particle accelerators or naturally by fast electrons moving through magnetic fields. The radiation produced in this way has a characteristic polarization, and the frequencies generated can range over a large portion of the electromagnetic spectrum.
Synchrotron radiation is similar to bremsstrahlung radiation, which is emitted by a charged particle when the acceleration is parallel to the direction of motion. The general term for radiation emitted by particles in a magnetic field is gyromagnetic radiation, for which synchrotron radiation is the ultra-relativistic special case. Radiation emitted by charged particles moving non-relativistically in a magnetic field is called cyclotron emission. For particles in the mildly relativistic range (≈85% of the speed of light), the emission is termed gyro-synchrotron radiation.
In astrophysics, synchrotron emission occurs, for instance, due to ultra-relativistic motion of a charged particle around a black hole. When the source follows a circular geodesic around the black hole, the synchrotron radiation occurs for orbits close to the photosphere where the motion is in the ultra-relativistic regime.
History
Synchrotron radiation was first observed by technician Floyd Haber, on April 24, 1947, at the 70 MeV electron synchrotron of the General Electric research laboratory in Schenectady, New York. While this was not the first synchrotron built, it was the first with a transparent vacuum tube, allowing the radiation to be directly observed.
As recounted by Herbert Pollock:
Description
A direct consequence of Maxwell's equations is that accelerated charged particles always emit electromagnetic radiation. Synchrotron radiation is the special case of charged particles moving at relativistic speed undergoing acceleration perpendicular to their direction of motion, typically in a magnetic field. In such a field, the force due to the field is always perpendicular to both the direction of motion and to the direction of field, as shown by the Lorentz force law.
The power carried by the radiation is found (in SI units) by the relativistic Larmor formula:
where
is the vacuum permittivity,
is the particle charge,
is the magnitude of the acceleration,
is the speed of light,
is the Lorentz factor,
,
is the radius of curvature of the particle trajectory.
The force on the emitting electron is given by the Abraham–Lorentz–Dirac force.
When the radiation is emitted by a particle moving in a plane, the radiation is linearly polarized when observed in that plane, and circularly polarized when observed at a small angle. Considering quantum mechanics, however, this radiation is emitted in discrete packets of photons and has significant effects in accelerators called quantum excitation. For a given acceleration, the average energy of emitted photons is proportional to and the emission rate to .
From accelerators
Circular accelerators will always produce gyromagnetic radiation as the particles are deflected in the magnetic field. However, the quantity and properties of the radiation are highly dependent on the nature of the acceleration taking place. For example, due to the difference in mass, the factor of in the formula for the emitted power means that electrons radiate energy at approximately 1013 times the rate of protons.
Energy loss from synchrotron radiation in circular accelerators was originally considered a nuisance, as additional energy must be supplied to the beam in order to offset the losses. However, beginning in the 1980s, circular electron accelerators known as light sources have been constructed to deliberately produce intense beams of synchrotron radiation for research.
In astronomy
Synchrotron radiation is also generated by astronomical objects, typically where relativistic electrons spiral (and hence change velocity) through magnetic fields.
Two of its characteristics include power-law energy spectra and polarization. It is considered to be one of the most powerful tools in the study of extra-solar magnetic fields wherever relativistic charged particles are present. Most known cosmic radio sources emit synchrotron radiation. It is often used to estimate the strength of large cosmic magnetic fields as well as analyze the contents of the interstellar and intergalactic media.
History of detection
This type of radiation was first detected in the Crab Nebula in 1956 by Jan Hendrik Oort and Theodore Walraven, and a few months later in a jet emitted by Messier 87 by Geoffrey R. Burbidge. It was confirmation of a prediction by Iosif S. Shklovsky in 1953. However, it had been predicted earlier (1950) by Hannes Alfvén and Nicolai Herlofson. Solar flares accelerate particles that emit in this way, as suggested by R. Giovanelli in 1948 and described by J.H. Piddington in 1952.
T. K. Breus noted that questions of priority on the history of astrophysical synchrotron radiation are complicated, writing:
From supermassive black holes
It has been suggested that supermassive black holes produce synchrotron radiation in "jets", generated by the gravitational acceleration of ions in their polar magnetic fields. The nearest such observed jet is from the core of the galaxy Messier 87. This jet is interesting for producing the illusion of superluminal motion as observed from the frame of Earth. This phenomenon is caused because the jets are traveling very near the speed of light and at a very small angle towards the observer. Because at every point of their path the high-velocity jets are emitting light, the light they emit does not approach the observer much more quickly than the jet itself. Light emitted over hundreds of years of travel thus arrives at the observer over a much smaller time period, giving the illusion of faster than light travel, despite the fact that there is actually no violation of special relativity.
Pulsar wind nebulae
A class of astronomical sources where synchrotron emission is important is pulsar wind nebulae, also known as plerions, of which the Crab nebula and its associated pulsar are archetypal.
Pulsed emission gamma-ray radiation from the Crab has recently been observed up to ≥25 GeV, probably due to synchrotron emission by electrons trapped in the strong magnetic field around the pulsar.
Polarization in the Crab nebula at energies from 0.1 to 1.0 MeV, illustrates this typical property of synchrotron radiation.
Interstellar and intergalactic media
Much of what is known about the magnetic environment of the interstellar medium and intergalactic medium is derived from observations of synchrotron radiation. Cosmic ray electrons moving through the medium interact with relativistic plasma and emit synchrotron radiation which is detected on Earth. The properties of the radiation allow astronomers to make inferences about the magnetic field strength and orientation in these regions. However, accurate calculations of field strength cannot be made without knowing the relativistic electron density.
In supernovae
When a star explodes in a supernova, the fastest ejecta move at semi-relativistic speeds approximately 10% the speed of light. This blast wave gyrates electrons in ambient magnetic fields and generates synchrotron emission, revealing the radius of the blast wave at the location of the emission. Synchrotron emission can also reveal the strength of the magnetic field at the front of the shock wave, as well as the circumstellar density it encounters, but strongly depends on the choice of energy partition between the magnetic field, proton kinetic energy, and electron kinetic energy. Radio synchrotron emission has allowed astronomers to shed light on mass loss and stellar winds that occur just prior to stellar death.
| Physical sciences | Electromagnetic radiation | Physics |
170825 | https://en.wikipedia.org/wiki/Shipworm | Shipworm | The shipworms, also called Teredo worms or simply Teredo (, via Latin ), are marine bivalve molluscs in the family Teredinidae, a group of saltwater clams with long, soft, naked bodies. They are notorious for boring into (and commonly eventually destroying) wood that is immersed in seawater, including such structures as wooden piers, docks, and ships; they drill passages by means of a pair of very small shells ("valves") borne at one end, with which they rasp their way through. They are sometimes called "termites of the sea". Carl Linnaeus assigned the common name Teredo to the best-known genus of shipworms in the 10th edition of his taxonomic magnum opus, Systema Naturæ (1758).
Characteristics
Removed from its burrow, the fully grown teredo ranges from several centimeters to about a meter in length, depending on the species. An average adult shipworm measures in length and less than in diameter, but some species grow to considerable size. The body is cylindrical, slender, naked, and superficially vermiform (worm-shaped). In spite of their slender, worm-like forms, shipworms possess the characteristic morphology of bivalves. The ctinidia lie mainly within the branchial siphon, through which the animal pumps the water that passes over the gills.
The two siphons are very long and protrude from the posterior end of the animal. Where they leave the end of the main part of the body, the siphons pass between a pair of calcareous plates called pallets. If the animal is alarmed, it withdraws the siphons and the pallets protectively block the opening of the tunnel.
The pallets are not to be confused with the two valves of the main shell, which are at the anterior end of the animal. Because they are the organs that the animal applies to boring its tunnel, they generally are located at the tunnel's end. They are borne on the slightly thickened, muscular anterior end of the cylindrical body and they are roughly triangular in shape and markedly concave on their interior surfaces. The outer surfaces are convex and in most species are deeply sculpted into sharp grinding surfaces with which the animals bore their way through the wood or similar medium in which they live and feed. The valves of shipworms are separated and the aperture of the mantle lies between them. The small "foot" (corresponding to the foot of a clam) can protrude through the aperture.
When shipworms bore into submerged wood, bacterial symbionts embedded within a sub-organ called the typhlosole in the shipworm gut, aid in the digestion of the wood particles ingested, The Alteromonas or Alteromonas-sub-group of bacteria identified as the symbiont species in the typhlosole, are known to digest lignin, and wood material in general. The tough molecular layers of lignin surround the cellulose elementary fibrils in the wood particles, and the lignin must be digested initially to allow access by other enzymes into the cellulose for digestion. Another bacterial species (Teredinibacter turnerae), in the gills secrete a variety of cellulose-digesting enzymes which may be secreted into the shipworm gut via a special organ called the gland of Deshayes. These secretions aid the shipworm's own carbohydrate-active enzymes (CAZymes) in digesting the wood particles in combination with the enzymes and potentially other metabolites secreted by the symbiont bacterial in the typhlosole.
The excavated burrow is usually lined with a calcareous tube. The valves of the shell of shipworms are small separate parts located at the anterior end of the worm, used for excavating the burrow. The protective role of the shells is lost because the animal spends all its life surrounded by wood.
Teredo navalis develops from eggs to metamorphosing larvae in about five weeks. They spend half of this time in the mother's gill chamber before being discharged as free-swimming larvae into the sea. Their sexes alternate, young are hermaphrodites while adults can be either male or female. Typically, organisms are male at first and female subsequently. A second male to female phase may occur, however shipworms rarely live long enough to complete the second phase. They have a lifespan of 1 to 3 years.
Anatomy
Shipworm anatomy reveals the typical organs of a bivalve mollusk, although with dimensional or positional peculiarities due to the thinness and length of the occupied space. Furthermore, some structures find no equivalent in other bivalve groups.
Gills are divided in two halves, the anterior one of small size, the posterior one much more developed. They are linked by the alimentary tract running on the side of the visceral mass.
The heart-kidney system is tilted, bringing the kidneys in a dorsal position relative to the heart, whose atria find themselves behind the ventricle. Furthermore, the anterior and posterior aorta become respectively posterior and anterior.
The anus opens at the end of a long anal tube.
The digestive gland is divided into several parts, with separate orifices in the stomach.
A vast caecum is linked to the stomach.
The digestive tube bears a very peculiar structure, the gland of Deshayes, probably homologous to salivary glands, which link to the oesophagus and stretch to the dorsal side of the posterior part of the gills.
The orifice of the gallery bears pallets with their own musculature.
The siphon retractor muscles are inserted on the calcareous covering of the gallery, and not on the shell's valves which are much further out.
The anterior and posterior anterior muscles have an antagonistic action.
Normally, the shipworm's body fills the entire length of the gallery, but the anterior region can retract itself slightly with respect to the latter's extremity. Without the gills, the viscera only cover one-fourth of the total length and only their anterior part is partially covered by the shell.
Taxonomy
Shipworms are marine animals in the phylum Mollusca, order Bivalvia, family Teredinidae. They were included in the now obsolete order Eulamellibranchiata, in which many documents still place them.
Ruth Turner of Harvard University was the leading 20th century expert on the Teredinidae; she published a detailed monograph on the family, the 1966 volume A Survey and Illustrated Catalogue of the Teredinidae published by the Museum of Comparative Zoology. More recently, the endosymbionts that are found in the gills have been subject to study the bioconversion of cellulose for fuel energy research.
Shipworm species comprise several genera, of which Teredo is the most commonly mentioned. The best known species is Teredo navalis. Historically, Teredo concentrations in the Caribbean Sea have been substantially higher than in most other salt water bodies.
Genera within the family Teridinidae include:
Bactronophorus Tapparone-Canefri, 1877
Bankia Gray, 1842
Dicyathifer Iredale, 1932
Kuphus Guettard, 1770
Lithoredo Shipway, Distel & Rosenberg, 2019
Lyrodus Binney, 1870
Nausitoria Wright, 1884
Neoteredo Bartsch, 1920
Nototeredo Bartsch, 1923
Psiloteredo Bartsch, 1922
Spathoteredo Moll, 1928
Teredo Linnaeus, 1758
Teredora Bartsch, 1921
Teredothyra Bartsch, 1921
Uperotus Guettard, 1770
Zachsia Bulatoff & Rjabtschikoff, 1933
Species
The Teredo genus has approximately 20 species that live in wooden materials such as logs, pilings, ship, and practically any other submerged wooden construction from temperate to tropical ocean zones. The species is thought to be native to the Atlantic Ocean and was once known as the Atlantic shipworm, although its exact origin is unknown. The longest marine bivalve, Kuphus polythalamia, was found from a lagoon near Mindanao island in the southeastern part of the Philippines, which belongs to the same group of mussels and clams. The existence of huge mollusks was established for centuries and studied by the scientists, based on the shells they left behind that were the size of baseball bats (length , diameter ). The bivalve is a rare creature that spends its life inside an elephant tusk-like hard shell made of calcium carbonate. It has a protective cap over its head which it reabsorbs to burrow into the mud for food. The case of the shipworm is not just the home of the black slimy worm. Instead, it acts as the primary source of nourishment in a non-traditional way. K. polythalamia sifts mud and sediment with its gills. Most shipworms are relatively smaller and feed on rotten wood. This shipworm instead relies on a beneficial symbiotic bacteria living in its gills. The bacteria use the hydrogen sulfide for energy to produce organic compounds that in turn feed the shipworms, similar to the process of photosynthesis used by green plants to convert the carbon dioxide in the air into simple carbon compounds. Scientists found that K. polythalamia cooperates with different bacteria than other shipworms, which could be the reason why it evolved from consuming rotten wood to living on hydrogen sulfide in the mud. The internal organs of the shipworm have shrunk from lack of use over the course of its evolution. The scientists are planning to study the microbes found in the single gill of K. polythalamia to find a new possible antimicrobial substance.
Habitat
Teredo navalis are a cosmopolitan species that can be found both in the Atlantic and Pacific oceans. Since they occupy wooden flotsam and natural driftwood such as dead tree trunks, they are spread as the wood is carried by currents. They also travel inside the wooden-hulled vessels that help increase their spread worldwide. However, the origin of T. navalis remains uncertain due to the widespread usage of ships in global trade and the resulting spreading of shipworms.
During the free-living larva stage, the species colonizes new habitats and spreads. Larvae are extremely sensitive to the presence of wood and will take advantage of any opportunity to attach to and penetrate wooden structures. In the Baltic Sea, free-floating piles carved by shipworms can be observed floating hundreds of kilometers away from the original wooden structures. The limiting element for propagation is salinity, which must be greater than 8% for successful reproduction. Consequently, freshwater is deadly to these invertebrates. Reproduction occurs during warm summer months, and the larvae mature for production in just eight weeks. Each year, several generations can be produced. Their ideal temperature range is and therefore T. navalis can be found in temperate and tropical zones.
The shipworm lives in waters with oceanic salinity. Accordingly, it is rare in the brackish Baltic Sea, where wooden shipwrecks are preserved for much longer than in the oceans.
The range of various species has changed over time based on human activity. Many waters in developed countries that had been plagued by shipworms were cleared of them by pollution from the Industrial Revolution and the modern era; as environmental regulation led to cleaner waters, shipworms have returned. Climate change has also changed the range of species; some once found only in warmer and more salty waters like the Caribbean have established habitats in the Mediterranean.
Cultural impact
Shipworms greatly damage wooden hulls and marine piling, and have been the subject of much study to find methods to avoid their attacks. Copper sheathing was used on wooden ships in the latter 18th century and afterwards, as a method of preventing damage by teredo worms. The first historically documented use of copper sheathing was experiments held by the British Royal Navy with , which was coppered in 1761 and thoroughly inspected after a two-year cruise. In a letter from the Navy Board to the Admiralty dated 31 August 1763 it was written "that so long as copper plates can be kept upon the bottom, the planks will be thereby entirely secured from the effects of the worm."
In the Netherlands the shipworm caused a crisis in the 18th century by attacking the timber that faced the sea dike. After that the dikes had to be faced with stones. In 2009, Teredo caused several minor collapses along the Hudson River waterfront in Hoboken, New Jersey, due to damage to underwater pilings.
In the early 19th century, engineer Marc Brunel observed that the shipworm's valves simultaneously enabled it to tunnel through wood and protected it from being crushed by the swelling timber. With that idea, he designed the first tunnelling shield, a modular iron tunnelling framework which enabled workers to tunnel through the unstable riverbed beneath the Thames. The Thames Tunnel was the first successful large tunnel built under a navigable river.
Henry David Thoreau's poem "Though All the Fates" pays homage to "New England's worm" which, in the poem, infests the hull of "[t]he vessel, though her masts be firm". In time, no matter what the ship carries or where she sails, the shipworm "her hulk shall bore, / [a]nd sink her in the Indian seas". The hull of the ship wrecked by a whale, inspiring Moby Dick, had been weakened by shipworms. In the Norse Saga of Erik the Red, Bjarni Herjólfsson, said to be the first European to discover the Americas, had his ship drift into the Irish Sea where it was eaten up by shipworms. He allowed half the crew to escape in a smaller boat covered in seal tar, while he stayed behind to drown with his men.
Use as seafood
In Palawan and Aklan in the Philippines, the shipworm is called and is eaten as a delicacy. It is prepared as kinilaw—that is, raw (cleaned) but marinated with vinegar or lime juice, chopped chili peppers and onions, a process very similar to shrimp ceviche. Similarly, T. navalis can be found inside the dead and rotten trunk of mangroves in West Papua, Indonesia. To the locals, the Kamoro tribe, it is referred to as and is considered as a delicacy in daily meals. It can be eaten fresh and raw (cleaned) or cooked (cleaned and boiled) as well and usually marinated with lime juice and chili peppers. Since T. navalis are related to clams, mussels, and oysters, the taste of the flesh has been compared to a wide variety of foods, from milk to oysters. Similarly, the delicacy is harvested, sold, and eaten from those taken by local natives in the mangrove forests of West Papua and some part of Borneo Island, Indonesia, and the central coastal peninsular regions of Thailand near Ko Phra Thong.
T. navalis grow faster than any other bivalve because it does not require much energy to create its small shell. They can grow to be about long in just six months. Mussels and oysters, on the other hand, with their much bigger shells, can take up to two years to reach harvestable size.
| Biology and health sciences | Bivalvia | Animals |
170829 | https://en.wikipedia.org/wiki/Shock%20absorber | Shock absorber | A shock absorber or damper is a mechanical or hydraulic device designed to absorb and damp shock impulses. It does this by converting the kinetic energy of the shock into another form of energy (typically heat) which is then dissipated. Most shock absorbers are a form of dashpot (a damper which resists motion via viscous friction).
Description
Pneumatic and hydraulic shock absorbers are used in conjunction with cushions and springs. An automobile shock absorber contains spring-loaded check valves and orifices to control the flow of oil through an internal piston (see below).
One design consideration, when designing or choosing a shock absorber, is where that energy will go. In most shock absorbers, energy is converted to heat inside the viscous fluid. In hydraulic cylinders, the hydraulic fluid heats up, while in air cylinders, the hot air is usually exhausted to the atmosphere. In other types of shock absorbers, such as electromagnetic types, the dissipated energy can be stored and used later. In general terms, shock absorbers help cushion vehicles on uneven roads and keep wheels in contact with the ground.
Vehicle suspension
In a vehicle, shock absorbers reduce the effect of traveling over rough ground, leading to improved ride quality and vehicle handling. While shock absorbers serve the purpose of limiting excessive suspension movement, their intended main purpose is to damp spring oscillations. Shock absorbers use valving of oil and gasses to absorb excess energy from the springs. Spring rates are chosen by the manufacturer based on the weight of the vehicle, loaded and unloaded. Some people use shocks to modify spring rates but this is not the correct use. Along with hysteresis in the tire itself, they damp the energy stored in the motion of the unsprung weight up and down. Effective wheel bounce damping may require tuning shocks to an optimal resistance.
Spring-based shock absorbers commonly use coil springs or leaf springs, though torsion bars are used in torsional shocks as well. Ideal springs alone, however, are not shock absorbers, as springs only store and do not dissipate or absorb energy. Vehicles typically employ both hydraulic shock absorbers and springs or torsion bars. In this combination, "shock absorber" refers specifically to the hydraulic piston that absorbs and dissipates vibration. Now, composite suspension systems are used mainly in 2 wheelers and also leaf springs are made up of composite material in 4 wheelers.
Construction
Shock absorbers are an important part of car suspension designed to increase comfort, stability and overall safety. The shock absorber, produced with precision and engineering skills, has many important features. The most common type is a hydraulic shock absorber, which usually includes a piston, a cylinder, and an oil-filled chamber. The piston is connected to the piston rod, which extends into the cylinder and divides the cylinder into two parts. One chamber is filled with hydraulic oil, while the other chamber contains or air. When there is an accident or vibration in the vehicle, the piston moves into the cylinder, forcing the hydraulic fluid through small holes, creating resistance and dissipating energy in the form of heat. This dampens oscillations, reducing further bouncing or wobble of the car. Shock construction requires a balance of features such as piston design, fluid viscosity, and overall size of the unit to ensure performance. As technology developed, other types of shock absorbers emerged, including gas and electric shock absorbers, that provided improved control and flexibility. The design and manufacture of shock absorbers is constantly evolving due to the continuous improvement of vehicle dynamics and passenger comfort.
Early history
In common with carriages and railway locomotives, most early motor vehicles used leaf springs. One of the features of these springs was that the friction between the leaves offered a degree of damping, and in a 1912 review of vehicle suspension, the lack of this characteristic in helical springs was the reason it was "impossible" to use them as main springs. However the amount of damping provided by leaf spring friction was limited and variable according to the conditions of the springs, and whether wet or dry. It also operated in both directions. Motorcycle front suspension adopted coil sprung Druid forks from about 1906, and similar designs later added Friction disk shock absorber rotary friction dampers, which damped both ways - but they were adjustable (e.g. 1924 Webb forks). These friction disk shock absorber s was also fitted to many cars.
One of the problems with motor cars was the large variation in sprung weight between lightly loaded and fully loaded, especially for the rear springs. When heavily loaded the springs could bottom out, and apart from fitting rubber 'bump stops', there were attempts to use heavy main springs with auxiliary springs to smooth the ride when lightly loaded, which were often called 'shock absorbers'. Realizing that the spring and vehicle combination bounced with a characteristic frequency, these auxiliary springs were designed with a different period, but were not a solution to the problem that the spring rebound after striking a bump could throw you out of your seat. What was called for was damping that operated on the rebound.
Although C.L. Horock came up with a design in 1901 that had hydraulic damping, it worked in one direction only. It does not seem to have gone into production right away, whereas mechanical dampers such as the Gabriel Snubber started being fitted in the late 1900s (also the similar Stromberg Anti-Shox). These used a belt coiled inside a device such that it freely wound in under the action of a coiled spring but met friction when drawn out. Gabriel Snubbers were fitted to an 11.9HP Arrol-Johnston car which broke the 6 hour Class B record at Brooklands in late 1912, and the Automator journal noted that this snubber might have a great future for racing due to its light weight and easy fitment.
French engineers Gaston Dumond and Ernest Mathis patented two different hydraulic shock absorbers with rectilinear motion in 1906–1907, but those were not commercially successful.
One of the earliest hydraulic dampers to go into production was the Telesco Shock Absorber, exhibited at the 1912 Olympia Motor Show and marketed by Polyrhoe Carburettors Ltd. This contained a spring inside the telescopic unit like the pure spring type 'shock absorbers' mentioned above, but also oil and an internal valve so that the oil damped in the rebound direction. The Telesco unit was fitted at the rear end of the leaf spring, in place of the rear spring to chassis mount, so that it formed part of the springing system, albeit a hydraulically damped part. This layout was presumably selected as it was easy to apply to existing vehicles, but it meant the hydraulic damping was not applied to the action of the main leaf spring, but only to the action of the auxiliary spring in the unit itself.
The first production hydraulic dampers to act on the main leaf spring movement were probably those based on an original concept by Maurice Houdaille patented in 1908 and 1909. These used a lever arm which moved hydraulically damped vanes inside the unit. The main advantage over the friction disk dampers was that it would resist sudden movement but allow slow movement, whereas the rotary friction dampers tended to stick and then offer the same resistance regardless of speed of movement. There appears to have been little progress on commercialising the lever arm shock absorbers until after World War I, after which they came into widespread use, for example as standard equipment on the 1927 Ford Model A and manufactured by Houde Engineering Corporation of Buffalo, NY.
Types of vehicle shock absorbers
Most vehicular shock absorbers are either twin-tube or mono-tube types with some variations on these themes.
Twin-tube
Basic twin-tube
Also known as a "two-tube" shock absorber, this device consists of two nested cylindrical tubes, an inner tube that is called the "working tube" or the "pressure tube", and an outer tube called the "reserve tube". At the bottom of the device on the inside is a compression valve or base valve. When the piston is forced up or down by bumps in the road, hydraulic fluid moves between different chambers via small holes or "orifices" in the piston and via the valve, converting the "shock" energy into heat which must then be dissipated.
Twin-tube gas charged
Variously known as a "gas cell two-tube" or similarly named design, this variation represented a significant advancement over the basic twin-tube form. Its overall structure is very similar to the twin-tube, but a low-pressure charge of nitrogen gas is added to the reserve tube. The result of this alteration is a dramatic reduction in "foaming" or "aeration", the undesirable outcome of a twin-tube overheating and failing which presents as foaming hydraulic fluid dripping out of the assembly. Twin-tube gas charged shock absorbers represent the vast majority of original modern vehicle suspension installations.
Position sensitive damping
Often abbreviated simply as "PSD", this design is another evolution of the twin-tube shock. In a PSD shock absorber, which still consists of two nested tubes and still contains nitrogen gas, a set of grooves has been added to the pressure tube. These grooves allow the piston to move relatively freely in the middle range of travel (i.e., the most common street or highway use, called by engineers the "comfort zone") and to move with significantly less freedom in response to shifts to more irregular surfaces when upward and downward movement of the piston starts to occur with greater intensity (i.e., on bumpy sections of roads— the stiffening gives the driver greater control of movement over the vehicle so its range on either side of the comfort zone is called the "control zone"). This advance allowed car designers to make a shock absorber tailored to specific makes and models of vehicles and to take into account a given vehicle's size and weight, its maneuverability, its horsepower, etc. in creating a correspondingly effective shock.
Acceleration sensitive damping
The next phase in shock absorber evolution was the development of a shock absorber that could sense and respond to not just situational changes from "bumpy" to "smooth" but to individual bumps in the road in a near instantaneous reaction. This was achieved through a change in the design of the compression valve, and has been termed "acceleration sensitive damping" or "ASD". Not only does this result in a complete disappearance of the "comfort vs. control" tradeoff, it also reduced pitch during vehicle braking and roll during turns. However, ASD shocks are usually only available as aftermarket changes to a vehicle and are only available from a limited number of manufacturers.
Coilover
Coilover shock absorbers are usually a kind of twin-tube gas charged shock absorber inside the helical road spring. They are common on motorcycles and scooter rear suspensions, and widely used on front and rear suspensions in cars.
Mono-tube
The principal design alternative to the twin-tube form has been the mono-tube shock absorber which was considered a revolutionary advancement when it appeared in the 1950s. As its name implies, the mono-tube shock, which is also a gas-pressurized shock and also comes in a coilover format, consists of only one tube, the pressure tube, though it has two pistons. These pistons are called the working piston and the dividing or floating piston, and they move in relative synchrony inside the pressure tube in response to changes in road smoothness. The two pistons also completely separate the shock's fluid and gas components. The mono-tube shock absorber is consistently a much longer overall design than the twin-tubes, making it difficult to mount in passenger cars designed for twin-tube shocks. However, unlike the twin-tubes, the mono-tube shock can be mounted either way—it does not have any directionality. It also does not have a compression valve, whose role has been taken up by the dividing piston, and although it contains nitrogen gas, the gas in a mono-tube shock is under high pressure (260-360 p.s.i. or so) which can actually help it to support some of the vehicle's weight, something which no other shock absorber is designed to do.
Mercedes became the first auto manufacturer to install mono-tube shocks as standard equipment on some of their cars starting in 1958. They were manufactured by Bilstein, patented the design and first appeared in 1954s. Because the design was patented, no other manufacturer could use it until 1971 when the patent expired.
Spool valve
Spool valve dampers are characterized by the use of hollow cylindrical sleeves with machined-in oil passages as opposed to traditional conventional flexible discs or shims. Spool valving can be applied with monotube, twin-tube, or position-sensitive packaging, and is compatible with electronic control.
Primary among benefits cited in Multimatic’s 2010 patent filing is the elimination of performance ambiguity associated with flexible shims, resulting in mathematically predictable, repeatable, and robust pressure-flow characteristics.
Remote reservoir/piggy-back
An extra tube or container of oil connected to the oil compartment of the (main) shock via a flexible pipe (remote reservoir) or inflexible pipe (piggy-back shock). Increases the amount of oil a shock can carry without increasing its length or thickness.
Bypass shock
Allows each section of suspension travel to have an independent suspension tune. Bypass shock, double bypass shock, triple bypass shock etc. Triple bypass would have a separate set of suspension tuning controls for each of its three sections of suspension travel: initial travel, mid-travel, full-travel.
Theoretical approaches
There are several commonly used principles behind shock absorption:
Hysteresis of structural material, for example the compression of rubber disks, stretching of rubber bands and cords, bending of steel springs, or twisting of torsion bars. Hysteresis is the tendency for otherwise elastic materials to rebound with less force than was required to deform them. Simple vehicles with no separate shock absorbers are damped, to some extent, by the hysteresis of their springs and frames.
Dry friction as used in wheel brakes, by using disks (classically made of leather) at the pivot of a lever, with friction forced by springs. Used in early automobiles such as the Ford Model T, up through some British cars of the 1940s and on the French Citroën 2CV in the 1950s. Although now considered obsolete, an advantage of this system is its mechanical simplicity; the degree of damping can be easily adjusted by tightening or loosening the screw clamping the disks, and it can be easily rebuilt with simple hand tools. A disadvantage is that the damping force tends not to increase with the speed of the vertical motion.
Solid state, tapered chain shock absorbers, using one or more tapered, axial alignment(s) of granular spheres, typically made of metals such as nitinol, in a casing. ,
Fluid friction, for example the flow of fluid through a narrow orifice (hydraulics), constitutes the vast majority of automotive shock absorbers. This design first appeared on Mors racing cars in 1902. One advantage of this type is, by using special internal valving, the absorber may be made relatively soft to compression (allowing a soft response to a bump) and relatively stiff to extension, controlling "rebound", which is the vehicle response to energy stored in the springs; similarly, a series of valves controlled by springs can change the degree of stiffness according to the velocity of the impact or rebound. Specialized shock absorbers for racing purposes may allow the front end of a dragster to rise with minimal resistance under acceleration, then strongly resist letting it settle, thereby maintaining a desirable rearward weight distribution for enhanced traction.
Compression of a gas, for example pneumatic shock absorbers, which can act like springs as the air pressure is building to resist the force on it. Enclosed gas is compressible, so equipment is less subject to shock damage. This concept was first applied in series production on Citroën cars in 1954. Today, many shock absorbers are pressurized with compressed nitrogen, to reduce the tendency for the oil to cavitate under heavy use. This causes foaming which temporarily reduces the damping ability of the unit. In very heavy duty units used for racing or off-road use, there may even be a secondary cylinder connected to the shock absorber to act as a reservoir for the oil and pressurized gas. In aircraft landing gear, air shock absorbers may be combined with hydraulic damping to reduce bounce. Such struts are called oleo struts (combining oil and air) .
Inertial resistance to acceleration, the Citroën 2CV had shock absorbers that damp wheel bounce with no external moving parts. These consisted of a spring-mounted 3.5 kg (7.75 lb) iron weight inside a vertical cylinder and are similar to, yet much smaller than versions of the tuned mass dampers used on tall buildings.
Composite hydropneumatic suspension combines many suspension elements in a single device: spring action, shock absorption, ride-height control, and self leveling suspension. This combines the advantages of gas compressibility and the ability of hydraulic machinery to apply force multiplication.
Conventional shock absorbers can be combined with air suspension springs - an alternate way to achieve ride-height control, and self leveling suspension.
In an electrorheological fluid damper, an electric field changes the viscosity of the oil. This principle allows semi-active damper applications in automotive and various industries.
Magnetic field variation: a magnetorheological damper changes its fluid characteristics through an electromagnet.
The effect of a shock absorber at high (sound) frequencies is usually limited by using a compressible gas as the working fluid or mounting it with rubber bushings.
Special features
Some shock absorbers allow tuning of the ride via control of the valve by a manual adjustment provided at the shock absorber.
In more expensive vehicles the valves may be remotely adjustable, offering the driver control of the ride at will while the vehicle is operated.
Additional control can be provided by dynamic valve control via computer in response to sensors, giving both a smooth ride and a firm suspension when needed, allowing ride height adjustment or even ride height control.
Ride height control is especially desirable in highway vehicles intended for occasional rough road use, as a means of improving handling and reducing aerodynamic drag by lowering the vehicle when operating on improved high speed roads.
Heatsinks, fans, or liquid cooling to prevent or delay shock fade and failure (oil leak) due to overheating
Shock absorber and strut comparison
A strut is a structural component that combines the shock absorber with other suspension parts like the coil spring and steering knuckle into one compact unit
Unlike a shock absorber, a strut has a reinforced body and stem.
Struts are subjected to multidirectional loads, while a shock absorber only damps vibration, only receiving a load along its axis.
Struts and shock absorbers have a different way of attachment. Shock absorbers are mounted through rubber or urethane bushings to the frame and suspension. A strut is hard mounted to the suspension and is mounted to the frame through a rotating plate providing the upper pivot point of the steering.
| Technology | Mechanisms | null |
170853 | https://en.wikipedia.org/wiki/Supergiant | Supergiant | Supergiants are among the most massive and most luminous stars. Supergiant stars occupy the top region of the Hertzsprung–Russell diagram with absolute visual magnitudes between about −3 and −8. The temperature range of supergiant stars spans from about 3,400 K to over 20,000 K.
Definition
The title supergiant, as applied to a star, does not have a single concrete definition. The term giant star was first coined by Hertzsprung when it became apparent that the majority of stars fell into two distinct regions of the Hertzsprung–Russell diagram. One region contained larger and more luminous stars of spectral types A to M and received the name giant. Subsequently, as they lacked any measurable parallax, it became apparent that some of these stars were significantly larger and more luminous than the bulk, and the term super-giant arose, quickly adopted as supergiant.
Supergiants with spectral classes of O to A are typically referred to as blue supergiants, supergiants with spectral classes F and G are referred to as yellow supergiants, while those of spectral classes K to M are red supergiants. Another convention uses temperature: supergiants with effective temperatures below 4800 K are deemed red supergiants; those with temperatures between 4800 and 7500 K are yellow supergiants, and those with temperatures exceeding 7500 K are blue supergiants. These correspond approximately to spectral types M and K for red supergiants, G, F, and late A for yellow supergiants, and early A, B, and O for blue supergiants.
Spectral luminosity class
Supergiant stars can be identified on the basis of their spectra, with distinctive lines sensitive to high luminosity and low surface gravity. In 1897, Antonia C. Maury had divided stars based on the widths of their spectral lines, with her class "c" identifying stars with the narrowest lines. Although it was not known at the time, these were the most luminous stars. In 1943, Morgan and Keenan formalised the definition of spectral luminosity classes, with class I referring to supergiant stars. The same system of MK luminosity classes is still used today, with refinements based on the increased resolution of modern spectra. Supergiants occur in every spectral class from young blue class O supergiants to highly evolved red class M supergiants. Because they are enlarged compared to main-sequence and giant stars of the same spectral type, they have lower surface gravities, and changes can be observed in their line profiles. Supergiants are also evolved stars with higher levels of heavy elements than main-sequence stars. This is the basis of the MK luminosity system which assigns stars to luminosity classes purely from observing their spectra.
In addition to the line changes due to low surface gravity and fusion products, the most luminous stars have high mass-loss rates and resulting clouds of expelled circumstellar materials which can produce emission lines, P Cygni profiles, or forbidden lines. The MK system assigns stars to luminosity classes: Ib for supergiants; Ia for luminous supergiants; and 0 (zero) or Ia+ for hypergiants. In reality there is much more of a continuum than well defined bands for these classifications, and classifications such as Iab are used for intermediate luminosity supergiants. Supergiant spectra are frequently annotated to indicate spectral peculiarities, for example B2 Iae or F5 Ipec.
Evolutionary supergiants
Supergiants can also be defined as a specific phase in the evolutionary history of certain stars. Stars with initial masses above quickly and smoothly initiate helium core fusion after they have exhausted their hydrogen, and continue fusing heavier elements after helium exhaustion until they develop an iron core, at which point the core collapses to produce a Type II supernova. Once these massive stars leave the main sequence, their atmospheres inflate, and they are described as supergiants. Stars initially under will never form an iron core and in evolutionary terms do not become supergiants, although they can reach luminosities thousands of times the sun's. They cannot fuse carbon and heavier elements after the helium is exhausted, so they eventually just lose their outer layers, leaving the core of a white dwarf. The phase where these stars have both hydrogen and helium burning shells is referred to as the asymptotic giant branch (AGB), as stars gradually become more and more luminous class M stars. Stars of may fuse sufficient carbon on the AGB to produce an oxygen-neon core and an electron-capture supernova, but astrophysicists categorise these as super-AGB stars rather than supergiants.
Categorisation of evolved stars
There are several categories of evolved stars that are not supergiants in evolutionary terms but may show supergiant spectral features or have luminosities comparable to supergiants.
Asymptotic-giant-branch (AGB) and post-AGB stars are highly evolved lower-mass red giants with luminosities that can be comparable to more massive red supergiants, but because of their low mass, being in a different stage of development (helium shell burning), and their lives ending in a different way (planetary nebula and white dwarf rather than supernova), astrophysicists prefer to keep them separate. The dividing line becomes blurred at around (or as high as in some models) where stars start to undergo limited fusion of elements heavier than helium. Specialists studying these stars often refer to them as super AGB stars, since they have many properties in common with AGB such as thermal pulsing. Others describe them as low-mass supergiants since they start to burn elements heavier than helium and can explode as supernovae. Many post-AGB stars receive spectral types with supergiant luminosity classes. For example, RV Tauri has an Ia (bright supergiant) luminosity class despite being less massive than the sun. Some AGB stars also receive a supergiant luminosity class, most notably W Virginis variables such as W Virginis itself, stars that are executing a blue loop triggered by thermal pulsing. A very small number of Mira variables and other late AGB stars have supergiant luminosity classes, for example α Herculis.
Classical Cepheid variables typically have supergiant luminosity classes, although only the most luminous and massive will actually go on to develop an iron core. The majority of them are intermediate mass stars fusing helium in their cores and will eventually transition to the asymptotic giant branch. δ Cephei itself is an example with a luminosity of and a mass of .
Wolf–Rayet stars are also high-mass luminous evolved stars, hotter than most supergiants and smaller, visually less bright but often more luminous because of their high temperatures. They have spectra dominated by helium and other heavier elements, usually showing little or no hydrogen, which is a clue to their nature as stars even more evolved than supergiants. Just as the AGB stars occur in almost the same region of the HR diagram as red supergiants, Wolf–Rayet stars can occur in the same region of the HR diagram as the hottest blue supergiants and main-sequence stars.
The most massive and luminous main-sequence stars are almost indistinguishable from the supergiants they quickly evolve into. They have almost identical temperatures and very similar luminosities, and only the most detailed analyses can distinguish the spectral features that show they have evolved away from the narrow early O-type main-sequence to the nearby area of early O-type supergiants. Such early O-type supergiants share many features with WNLh Wolf–Rayet stars and are sometimes designated as slash stars, intermediates between the two types.
Luminous blue variables (LBVs) stars occur in the same region of the HR diagram as blue supergiants but are generally classified separately. They are evolved, expanded, massive, and luminous stars, often hypergiants, but they have very specific spectral variability, which defies the assignment of a standard spectral type. LBVs observed only at a particular time or over a period of time when they are stable, may simply be designated as hot supergiants or as candidate LBVs due to their luminosity.
Hypergiants are frequently treated as a different category of star from supergiants, although in all important respects they are just a more luminous category of supergiant. They are evolved, expanded, massive and luminous stars like supergiants, but at the most massive and luminous extreme, and with particular additional properties of undergoing high mass-loss due to their extreme luminosities and instability. Generally only the more evolved supergiants show hypergiant properties, since their instability increases after high mass-loss and some increase in luminosity.
Some B[e] stars are supergiants although other B[e] stars are clearly not. Some researchers distinguish the B[e] objects as separate from supergiants, while researchers prefer to define massive evolved B[e] stars as a subgroup of supergiants. The latter has become more common with the understanding that the B[e] phenomenon arises separately in a number of distinct types of stars, including some that are clearly just a phase in the life of supergiants.
Properties
Supergiants have masses from 8 to 12 times the Sun () upwards, and luminosities from about 1,000 to over a million times the Sun (). They vary greatly in radius, usually from 30 to 500, or even in excess of 1,000 solar radii (). They are massive enough to begin helium-core burning gently before the core becomes degenerate, without a flash and without the strong dredge-ups that lower-mass stars experience. They go on to successively ignite heavier elements, usually all the way to iron. Also because of their high masses, they are destined to explode as supernovae.
The Stefan–Boltzmann law dictates that the relatively cool surfaces of red supergiants radiate much less energy per unit area than those of blue supergiants; thus, for a given luminosity, red supergiants are larger than their blue counterparts. Radiation pressure limits the largest cool supergiants to around 1,500 and the most massive hot supergiants to around a million (Mbol around −10). Stars near and occasionally beyond these limits become unstable, pulsate, and experience rapid mass loss.
Surface gravity
The supergiant luminosity class is assigned on the basis of spectral features that are largely a measure of surface gravity, although such stars are also affected by other properties such as microturbulence. Supergiants typically have surface gravities of around log(g) 2.0 cgs and lower, although bright giants (luminosity class II) have statistically very similar surface gravities to normal Ib supergiants. Cool luminous supergiants have lower surface gravities, with the most luminous (and unstable) stars having log(g) around zero. Hotter supergiants, even the most luminous, have surface gravities around one, due to their higher masses and smaller radii.
Temperature
There are supergiant stars at all of the main spectral classes and across the whole range of temperatures from mid-M class stars at around 3,400 K to the hottest O class stars over 40,000 K. Supergiants are generally not found cooler than mid-M class. This is expected theoretically since they would be catastrophically unstable; however, there are potential exceptions among extreme stars such as VX Sagittarii.
Although supergiants exist in every class from O to M, the majority are spectral type B (blue supergiants), more than at all other spectral classes combined. A much smaller grouping consists of very low-luminosity G-type supergiants, intermediate mass stars burning helium in their cores before reaching the asymptotic giant branch. A distinct grouping is made up of high-luminosity supergiants at early B (B0-2) and very late O (O9.5), more common even than main sequence stars of those spectral types. The number of post-main sequence blue supergiants is greater than those expected from theoretical models, leading to the "blue supergiant problem".
The relative numbers of blue, yellow, and red supergiants is an indicator of the speed of stellar evolution and is used as a powerful test of models of the evolution of massive stars.
Luminosity
The supergiants lie more or less on a horizontal band occupying the entire upper portion of the HR diagram, but there are some variations at different spectral types. These variations are due partly to different methods for assigning luminosity classes at different spectral types, and partly to actual physical differences in the stars.
The bolometric luminosity of a star reflects its total output of electromagnetic radiation at all wavelengths. For very hot and very cool stars, the bolometric luminosity is dramatically higher than the visual luminosity, sometimes several magnitudes or a factor of five or more. This bolometric correction is approximately one magnitude for mid B, late K, and early M stars, increasing to three magnitudes (a factor of 15) for O and mid M stars.
All supergiants are larger and more luminous than main sequence stars of the same temperature. This means that hot supergiants lie on a relatively narrow band above bright main sequence stars. A B0 main sequence star has an absolute magnitude of about −5, meaning that all B0 supergiants are significantly brighter than absolute magnitude −5. Bolometric luminosities for even the faintest blue supergiants are tens of thousands of times the sun (). The brightest can be and are often unstable such as α Cygni variables and luminous blue variables.
The very hottest supergiants with early O spectral types occur in an extremely narrow range of luminosities above the highly luminous early O main sequence and giant stars. They are not classified separately into normal (Ib) and luminous (Ia) supergiants, although they commonly have other spectral type modifiers such as "f" for nitrogen and helium emission (e.g. O2 If for HD 93129A).
Yellow supergiants can be considerably fainter than absolute magnitude −5, with some examples around −2 (e.g. 14 Persei). With bolometric corrections around zero, they may only be a few hundred times the luminosity of the sun. These are not massive stars, though; instead, they are stars of intermediate mass that have particularly low surface gravities, often due to instability such as Cepheid pulsations. These intermediate mass stars' being classified as supergiants during a relatively long-lasting phase of their evolution account for the large number of low luminosity yellow supergiants. The most luminous yellow stars, the yellow hypergiants, are amongst the visually brightest stars, with absolute magnitudes around −9, although still less than .
There is a strong upper limit to the luminosity of red supergiants at around . Stars that would be brighter than this shed their outer layers so rapidly that they remain hot supergiants after they leave the main sequence. The majority of red supergiants were main sequence stars and now have luminosities below , and there are very few bright supergiant (Ia) M class stars. The least luminous stars classified as red supergiants are some of the brightest AGB and post-AGB stars, highly expanded and unstable low mass stars such as the RV Tauri variables. The majority of AGB stars are given giant or bright giant luminosity classes, but particularly unstable stars such as W Virginis variables may be given a supergiant classification (e.g. W Virginis itself). The faintest red supergiants are around absolute magnitude −3.
Variability
While most supergiants such as Alpha Cygni variables, semiregular variables, and irregular variables show some degree of photometric variability, certain types of variables amongst the supergiants are well defined. The instability strip crosses the region of supergiants, and specifically many yellow supergiants are Classical Cepheid variables. The same region of instability extends to include the even more luminous yellow hypergiants, an extremely rare and short-lived class of luminous supergiant. Many R Coronae Borealis variables, although not all, are yellow supergiants, but this variability is due to their unusual chemical composition rather than a physical instability.
Further types of variable stars such as RV Tauri variables and PV Telescopii variables are often described as supergiants. RV Tau stars are frequently assigned spectral types with a supergiant luminosity class on account of their low surface gravity, and they are amongst the most luminous of the AGB and post-AGB stars, having masses similar to the sun; likewise, the even rarer PV Tel variables are often classified as supergiants, but have lower luminosities than supergiants and peculiar B[e] spectra extremely deficient in hydrogen. Possibly they are also post-AGB objects or "born-again" AGB stars.
The LBVs are variable with multiple semi-regular periods and less predictable eruptions and giant outbursts. They are usually supergiants or hypergiants, occasionally with Wolf-Rayet spectra—extremely luminous, massive, evolved stars with expanded outer layers, but they are so distinctive and unusual that they are often treated as a separate category without being referred to as supergiants or given a supergiant spectral type. Often their spectral type will be given just as "LBV" because they have peculiar and highly variable spectral features, with temperatures varying from about 8,000 K in outburst up to 20,000 K or more when "quiescent."
Chemical abundances
The abundance of various elements at the surface of supergiants is different from less luminous stars. Supergiants are evolved stars and may have undergone convection of fusion products to the surface.
Cool supergiants show enhanced helium and nitrogen at the surface due to convection of these fusion products to the surface during the main sequence of very massive stars, to dredge-ups during shell burning, and to the loss of the outer layers of the star. Helium is formed in the core and shell by fusion of hydrogen and nitrogen which accumulates relative to carbon and oxygen during CNO cycle fusion. At the same time, carbon and oxygen abundances are reduced. Red supergiants can be distinguished from luminous but less massive AGB stars by unusual chemicals at the surface, enhancement of carbon from deep third dredge-ups, as well as carbon-13, lithium and s-process elements. Late-phase AGB stars can become highly oxygen-enriched, producing OH masers.
Hotter supergiants show differing levels of nitrogen enrichment. This may be due to different levels of mixing on the main sequence due to rotation or because some blue supergiants are newly evolved from the main sequence while others have previously been through a red supergiant phase. Post-red supergiant stars have a generally higher level of nitrogen relative to carbon due to convection of CNO-processed material to the surface and the complete loss of the outer layers. Surface enhancement of helium is also stronger in post-red supergiants, representing more than a third of the atmosphere.
Evolution
O type main-sequence stars and the most massive of the B type blue-white stars become supergiants. Due to their extreme masses, they have short lifespans, between 30 million years and a few hundred thousand years. They are mainly observed in young galactic structures such as open clusters, the arms of spiral galaxies, and in irregular galaxies. They are less abundant in spiral galaxy bulges and are rarely observed in elliptical galaxies, or globular clusters, which are composed mainly of old stars.
Supergiants develop when massive main-sequence stars run out of hydrogen in their cores, at which point they start to expand, just like lower-mass stars. Unlike lower-mass stars, however, they begin to fuse helium in the core smoothly and not long after exhausting their hydrogen. This means that they do not increase their luminosity as dramatically as lower-mass stars, and they progress nearly horizontally across the HR diagram to become red supergiants. Also unlike lower-mass stars, red supergiants are massive enough to fuse elements heavier than helium, so they do not puff off their atmospheres as planetary nebulae after a period of hydrogen and helium shell burning; instead, they continue to burn heavier elements in their cores until they collapse. They cannot lose enough mass to form a white dwarf, so they will leave behind a neutron star or black hole remnant, usually after a core collapse supernova explosion.
Stars more massive than about cannot expand into a red supergiant. Because they burn too quickly and lose their outer layers too quickly, they reach the blue supergiant stage, or perhaps yellow hypergiant, before returning to become hotter stars. The most massive stars, above about , hardly move at all from their position as O main-sequence stars. These convect so efficiently that they mix hydrogen from the surface right down to the core. They continue to fuse hydrogen until it is almost entirely depleted throughout the star, then rapidly evolve through a series of stages of similarly hot and luminous stars: supergiants, slash stars, WNh-, WN-, and possibly WC- or WO-type stars. They are expected to explode as supernovae, but it is not clear how far they evolve before this happens. The existence of these supergiants still burning hydrogen in their cores may necessitate a slightly more complex definition of supergiant: a massive star with increased size and luminosity due to fusion products building up, but still with some hydrogen remaining.
The first stars in the universe are thought to have been considerably brighter and more massive than the stars in the modern universe. Part of the theorized population III of stars, their existence is necessary to explain observations of elements other than hydrogen and helium in quasars. Possibly larger and more luminous than any supergiant known today, their structure was quite different, with reduced convection and less mass loss. Their very short lives are likely to have ended in violent photodisintegration or pair instability supernovae.
Supernova progenitors
Most type II supernova progenitors are thought to be red supergiants, while the less common type Ib/c supernovae are produced by hotter Wolf–Rayet stars that have completely lost more of their hydrogen atmosphere. Almost by definition, supergiants are destined to end their lives violently. Stars large enough to start fusing elements heavier than helium do not seem to have any way to lose enough mass to avoid catastrophic core collapse, although some may collapse, almost without trace, into their own central black holes.
The simple "onion" models showing red supergiants inevitably developing to an iron core and then exploding have been shown, however, to be too simplistic. The progenitor for the unusual type II Supernova 1987A was a blue supergiant, thought to have already passed through the red supergiant phase of its life, and this is now known to be far from an exceptional situation. Much research is now focused on how blue supergiants can explode as a supernova and when red supergiants can survive to become hotter supergiants again.
Well known examples
Supergiants are rare and short-lived stars, but their high luminosity means that there are many naked-eye examples, including some of the brightest stars in the sky. Rigel, the brightest star in the constellation Orion is a typical blue-white supergiant; the three stars of Orion's Belt are all blue supergiants; Deneb is the brightest star in Cygnus, another blue supergiant; and Delta Cephei (itself the prototype) and Polaris are Cepheid variables and yellow supergiants. Antares and VV Cephei A are red supergiants. μ Cephei is considered a red hypergiant due to its large luminosity and it is one of the reddest stars visible to the naked eye and one of the largest in the galaxy. Rho Cassiopeiae, a variable, yellow hypergiant, is one of the most luminous naked-eye stars. Betelgeuse is a red supergiant that may have been a yellow supergiant in antiquity and the second brightest star in the constellation Orion.
| Physical sciences | Stellar astronomy | null |
170920 | https://en.wikipedia.org/wiki/Shiitake | Shiitake | The shiitake (; Chinese/black forest mushrooms or Lentinula edodes) is an edible mushroom native to East Asia, which is cultivated and consumed around the globe.
Taxonomy
The fungus was first described scientifically as Agaricus edodes by Miles Joseph Berkeley in 1877. It was placed in the genus Lentinula by David Pegler in 1976. The fungus has acquired an extensive synonymy in its taxonomic history:
Agaricus edodes Berk. (1878)
Armillaria edodes (Berk.) Sacc. (1887)
Mastoleucomychelloes edodes (Berk.) Kuntze (1891)
Cortinellus edodes (Berk.) S.Ito & S.Imai (1938)
Lentinus edodes (Berk.) Singer (1941)
Collybia shiitake J.Schröt. (1886)
Lepiota shiitake (J.Schröt.) Nobuj. Tanaka (1889)
Cortinellus shiitake (J.Schröt.) Henn. (1899)
Tricholoma shiitake (J.Schröt.) Lloyd (1918)
Lentinus shiitake (J.Schröt.) Singer (1936)
Lentinus tonkinensis Pat. (1890)
Lentinus mellianus Lohwag (1918)
The mushroom's Japanese name is a compound word composed of , for the tree Castanopsis cuspidata that provides the dead logs on which it is typically cultivated, and . The specific epithet is the Latin word for "edible".
It is also commonly called "sawtooth oak mushroom", "black forest mushroom", "black mushroom", "golden oak mushroom", or "oakwood mushroom".
Distribution and habitat
Shiitake grow in groups on the decaying wood of deciduous trees, particularly shii and other chinquapins, chestnut, oak, maple, beech, sweetgum, poplar, hornbeam, ironwood, and mulberry. Its natural distribution includes warm and moist climates in Southeast Asia.
Cultivation
The earliest written record of shiitake cultivation is seen in the Records of Longquan County () compiled by He Zhan () in 1209 during the Song dynasty in China. The 185-word description of shiitake cultivation from that literature was later cross-referenced many times and eventually adapted in a book by a Japanese horticulturist in 1796, the first book on shiitake cultivation in Japan.
The Japanese cultivated the mushroom by cutting shii trees with axes and placing the logs by trees that were already growing shiitake or contained shiitake spores. Before 1982, the Japan Islands' variety of these mushrooms could only be grown in traditional locations using ancient methods. A 1982 report on the budding and growth of the Japanese variety revealed opportunities for commercial cultivation in the United States.
Shiitake are widely cultivated worldwide, contributing about 25% of the total yearly production of mushrooms. Commercially, shiitake mushrooms are typically grown in conditions similar to their natural environment on either artificial substrate or hardwood logs, such as oak.
Toxicity
Rarely, consumption of raw or slightly cooked shiitake mushrooms may cause an allergic reaction called "shiitake dermatitis", including an erythematous, micro-papular, streaky pruriginous rash that occurs all over the body including face and scalp, appearing about 24 hours after consumption, possibly worsening by sun exposure and disappearing after 3 to 21 days. This effect – presumably caused by the polysaccharide, lentinan – is more common in East Asia, but may be growing in occurrence in Europe as shiitake consumption increases. Thorough cooking may eliminate the allergenicity.
Uses
Fresh and dried shiitake have many uses in East and Southeast Asia. In Chinese cuisine, they are used in many dishes, including soups, braises, and stir-fried vegetable dishes such as Buddha's delight. In Japan, they are served in miso soup, used as the basis for a kind of vegetarian dashi, and as an ingredient in many steamed and simmered dishes.
Two prized varieties are produced in cooler temperatures: One high-grade variety is called () (literally "winter mushroom") in Chinese, or in Japanese. The most highly prized variety is called () (literally "flower mushroom") in Chinese, due to the flower-like pattern of cracks in the cap.
Nutrition
In a reference serving, raw shiitake mushrooms provide of food energy and are 90% water, 7% carbohydrates, 2% protein and less than 1% fat. Raw shiitake mushrooms contain moderate levels of some dietary minerals.
Like all mushrooms, shiitakes produce vitamin D2 upon exposure of their internal ergosterol to ultraviolet B (UVB) rays from sunlight or broadband UVB fluorescent tubes.
Gallery
| Biology and health sciences | Edible fungi | null |
170927 | https://en.wikipedia.org/wiki/Whooping%20cough | Whooping cough | Whooping cough ( or ), also known as pertussis or the 100-day cough, is a highly contagious, vaccine-preventable bacterial disease. Initial symptoms are usually similar to those of the common cold with a runny nose, fever, and mild cough, but these are followed by two or three months of severe coughing fits. Following a fit of coughing, a high-pitched whoop sound or gasp may occur as the person breathes in. The violent coughing may last for 10 or more weeks, hence the phrase "100-day cough". The cough may be so hard that it causes vomiting, rib fractures, and fatigue. Children less than one year old may have little or no cough and instead have periods when they cannot breathe. The incubation period is usually seven to ten days. Disease may occur in those who have been vaccinated, but symptoms are typically milder.
The bacterium Bordetella pertussis causes pertussis, which is spread easily through the coughs and sneezes of an infected person. People are infectious from the start of symptoms until about three weeks into the coughing fits. Diagnosis is by collecting a sample from the back of the nose and throat. This sample can then be tested either by culture or by polymerase chain reaction.
Prevention is mainly by vaccination with the pertussis vaccine. Initial immunization is recommended between six and eight weeks of age, with four doses to be given in the first two years of life. Protection from pertussis decreases over time, so additional doses of vaccine are often recommended for older children and adults. Vaccination during pregnancy is highly effective at protecting the infant from pertussis during their vulnerable early months of life, and is recommended in many countries. Antibiotics may be used to prevent the disease in those who have been exposed and are at risk of severe disease. In those with the disease, antibiotics are useful if started within three weeks of the initial symptoms, but otherwise have little effect in most people. In pregnant women and children less than one year old, antibiotics are recommended within six weeks of symptom onset. Antibiotics used include erythromycin, azithromycin, clarithromycin, or trimethoprim/sulfamethoxazole. Evidence to support interventions for the cough, other than antibiotics, is poor. About 50% of infected children less than a year old require hospitalization and nearly 0.5% (1 in 200) die.
An estimated 16.3 million people worldwide were infected in 2015. Most cases occur in the developing world, and people of all ages may be affected. In 2015, pertussis resulted in 58,700 deaths – down from 138,000 deaths in 1990. Outbreaks of the disease were first described in the 16th century. The bacterium that causes the infection was discovered in 1906. The pertussis vaccine became available in the 1940s.
Signs and symptoms
The classic symptoms of pertussis are a paroxysmal cough, inspiratory whoop, and fainting, or vomiting after coughing. The cough from pertussis has been documented to cause subconjunctival hemorrhages, rib fractures, urinary incontinence, hernias, and vertebral artery dissection. Violent coughing can cause the pleura to rupture, leading to a pneumothorax. Vomiting after a coughing spell or an inspiratory whooping sound on coughing almost doubles the likelihood that the illness is pertussis. The absence of a paroxysmal cough or posttussive emesis makes it almost half as likely.
The illness usually starts with mild respiratory symptoms including mild coughing, sneezing, or a runny nose (known as the catarrhal stage). After one or two weeks, the coughing classically develops into uncontrollable fits, sometimes followed by a high-pitched "whoop" sound, as the person tries to inhale. About 50% of children and adults "whoop" at some point in diagnosed pertussis cases during the paroxysmal stage. This stage usually lasts up to 3 months, or sometimes longer. A gradual transition then occurs to the convalescent stage, which usually lasts one to four weeks. A decrease in paroxysms of coughing marks this stage, although paroxysms may occur with subsequent respiratory infection for many months after the onset of pertussis.
Symptoms of pertussis can be variable, especially between immunized and non-immunized people. Immunized people can present with a milder infection; they may only have the paroxysmal cough for a couple of weeks and may lack the "whooping" characteristic. Although immunized people have a milder form of the infection, they can still spread the disease to others who are not immune.
Incubation period
The time between exposure and the development of symptoms is on average 7–14 days (ranging 6–20 days), and rarely as long as 42 days.
Cause
Pertussis is caused by the bacterium Bordetella pertussis. It is an airborne disease (through droplets) that spreads easily through the coughs and sneezes of an infected person.
Host species
Humans are the only host species of B. pertussis. Outbreaks of whooping cough have been observed among chimpanzees in a zoo and wild gorillas; in both cases, it is considered likely that the infection was acquired as a result of close contact with humans. Several zoos have a long-standing custom of vaccinating their primates against whooping cough.
Mechanism
After the bacteria are inhaled, they initially adhere to the ciliated epithelium in the nasopharynx. Surface proteins of B. pertussis, including filamentous hemagglutinin and pertactin, mediate attachment to the epithelium. The bacteria then multiply. In infants, who experience more severe disease, the bacteria spread down to the lungs.
The bacteria secrete several toxins. Tracheal cytotoxin (TCT), a fragment of peptidoglycan, kills ciliated epithelial cells in the airway and thereby inhibits the mechanism which clears the airways of mucus and debris. TCT may contribute to the cough characteristic of pertussis. Pertussis toxin causes lymphocytosis by an unknown mechanism. The elevated number of white blood cells leads to pulmonary hypertension, a major cause of death by pertussis. In infants who develop encephalopathy, cerebral hemorrhage and cortical atrophy occur, likely due to hypoxia.
Diagnosis
Based on symptoms
A physician's overall impression is most effective in initially making the diagnosis. Single factors are much less useful. In adults with a cough of less than 8 weeks, vomiting after coughing or a "whoop" is supportive. If there are no bouts of coughing or there is a fever the diagnosis is unlikely. In children who have a cough of less than 4 weeks vomiting after coughing is somewhat supportive but not definitive.
Lab tests
Methods used in laboratory diagnosis include culturing of nasopharyngeal swabs on a nutrient medium (Bordet–Gengou medium), polymerase chain reaction (PCR), direct fluorescent antibody (DFA), and serological methods (e.g. complement fixation test). The bacteria can be recovered from the person only during the first three weeks of illness, rendering culturing and DFA useless after this period. However, PCR may have some limited usefulness for an additional three weeks.
Serology may be used for adults and adolescents who have already been infected for several weeks to determine whether antibodies against pertussis toxin or another virulence factor of B. pertussis are present at high levels in the person's blood.
Differential diagnosis
A similar, milder disease is caused by B. parapertussis.
Prevention
The primary method of prevention for pertussis is vaccination. Evidence is insufficient to determine the effectiveness of antibiotics in those who have been exposed, but are without symptoms. Preventive antibiotics, however, are still frequently used in those who have been exposed and are at high risk of severe disease (such as infants).
Vaccine
Pertussis vaccines are effective at preventing illness and are recommended for routine use by the World Health Organization and the United States Centers for Disease Control and Prevention. The vaccine saved an estimated half a million lives in 2002.
The multi-component acellular pertussis vaccine is 71–85% effective, with greater effectiveness against more severe strains. Despite widespread vaccination, pertussis has persisted in vaccinated populations. It remains "one of the most common vaccine-preventable diseases in Western countries". The 21st-century resurgence in pertussis infections is attributed to a combination of waning immunity and bacterial mutations that elude vaccines.
Immunization does not confer lifelong immunity; a 2011 CDC study indicated that protection may only last three to six years. This covers childhood, which is the time of greatest exposure and greatest risk of death from pertussis.
An effect of widespread immunization on society has been the shift of reported infections from children aged 1–9 years to infants, adolescents, and adults, with adolescents and adults acting as reservoirs for B. pertussis and infecting infants who have had fewer than three doses of vaccine.
Infection induces incomplete natural immunity that wanes over time. A 2005 study said estimates of the duration of infection-acquired immunity range from 7 to 20 years and the different results could be the result of differences in levels of circulating B. pertussis, surveillance systems, and case definitions used. The study said protective immunity after vaccination wanes after 4–12 years. One study suggested that the availability of vaccine exemptions increases the number of pertussis cases.
Some studies have suggested that while acellular pertussis vaccines effectively prevent disease, they have a limited impact on infection and transmission, meaning that vaccinated people could spread pertussis even though they may have only mild symptoms or none at all. Pertussis infection in these persons may be asymptomatic, or present as illness ranging from a mild cough to classic pertussis with persistent cough (i.e., lasting more than 7 days). Even though the disease may be milder in older persons, those who are infected may transmit the disease to other susceptible persons, including unimmunized or incompletely immunized infants. Older persons are often found to have the first case in a household with multiple pertussis cases and are often the source of infection for children.
Treatment
The antibiotics erythromycin, clarithromycin, or azithromycin are typically the recommended treatment. Newer macrolides are frequently recommended due to lower rates of side effects. Trimethoprim-sulfamethoxazole (TMP/SMX) may be used in those with allergies to first-line agents or in infants who have a risk of pyloric stenosis from macrolides.
A reasonable guideline is to treat people aged more than a year within three weeks of cough onset, infants aged less than one year, and pregnant women within six weeks of cough onset. If the person is diagnosed late, antibiotics will not alter the course of the illness, and even without antibiotics, they should no longer be spreading pertussis. When used early, antibiotics decrease the duration of infectiousness, and thus prevent spread. Short-term antibiotics (azithromycin for 3–5 days) are as effective as long-term treatment (erythromycin 10–14 days) in eliminating B. pertussis with fewer and less severe side effects.
People with pertussis are most infectious during the first two weeks following the onset of symptoms.
Effective treatments of the cough associated with this condition have not been developed. The use of over-the-counter cough medications is discouraged and has not been found helpful.
Prognosis
While most healthy older children and adults fully recover, infection in newborns is particularly severe. Pertussis is fatal in an estimated 0.5% of US infants under one year of age. First-year infants are also more likely to develop complications, such as apneas (31%), pneumonia (12%), seizures (0.6%) and encephalopathy (0.15%). This may be due to the ability of the bacterium to suppress the immune system.
Epidemiology
Pertussis is endemic worldwide. More than 151,000 cases were reported globally in 2018. However not all cases are reported or correctly diagnosed, especially in developing countries. Pertussis is one of the leading causes of vaccine-preventable deaths worldwide. A study in 2017 estimated the global burden of the disease to be 24 million cases per year with 160,000 deaths among young children, with about 90% of all cases occurring in developing countries.
Epidemics of the disease occur cyclically, every three to 5 years, both in areas with vaccination programs and those without. Over time, immunity declines in those who have either been vaccinated or have recovered from infection. In addition, infants are born who are susceptible to infection. An epidemic can occur once herd immunity decreases below a certain level. It is also possible that the bacterium is evolving to evade vaccine-induced immunity.
Before vaccines, an average of 178,171 cases was reported in the U.S., with peaks reported every two to five years; more than 93% of reported cases occurred in children under 10 years of age. With the widespread introduction of the DTP combined vaccine (diphtheria tetanus and pertussis) in the 1940s, pertussis incidence fell dramatically to approximately 1,000 by 1976, when they fluctuated between 1,000 and 30,000 annually.
Cases recorded outside of the U.S. were also recorded at high numbers comparable to their populations. Before the vaccine was discovered, Sweden averaged nearly 3,000 children deaths per year. With their population only being 1.8 million in the years 1749-64 this number was very high. The London population during the same period recorded over 3,000 deaths. The rates in London continued to grow into the 18th century. These numbers show how the disease affected not only the U.S. but also those around the world.
According to the CDC, reports that cases of whooping cough have reached their highest levels since 2014. This year, there have been over 16,000 cases, marking a fourfold increase compared to last year’s total of more than 3,700 cases. The CDC has also confirmed two deaths related to the illness. The United States is seeing a return to pre-pandemic trends, where annual cases typically exceed 10,000.
History
Discovery
B. pertussis was discovered in 1906 by Jules Bordet and Octave Gengou (the bacterium is subsequently named Bordetella pertussis in honour of Jules Bordet). They were able to successfully culture B. pertussis and went on to develop the first inactivated whole-cell vaccine in 1912, followed by other researchers in 1913 and 1914. These early vaccines had limited effectiveness. In the 1920s, Louis W. Sauer developed a vaccine for whooping cough at Evanston Hospital. In 1925 Danish physician Thorvald Madsen was the first to test a whole-cell vaccine on a wide scale. Madsen used the vaccine to control outbreaks in the Faroe Islands in the North Sea, however, two children died shortly after receiving the vaccine.
Vaccine
In 1932, an outbreak of whooping cough hit Atlanta, Georgia, prompting pediatrician Leila Denmark to begin her study of the disease. Over the next six years, her work was published in the Journal of the American Medical Association, and in partnership with Emory University and Eli Lilly & Company, she developed the first safe and effective pertussis vaccine. In 1942, American scientists Grace Eldering, Loney Gordon, and Pearl Kendrick combined the whole-cell pertussis vaccine with diphtheria and tetanus toxoids to generate the first DTP combination vaccine. To minimize the frequent side effects caused by the pertussis component, Japanese scientist Yuji Sato developed an acellular vaccine consisting of purified haemagglutinins (HAs: filamentous strep throat and leukocytosis-promoting-factor HA), which are secreted by B. pertussis. Sato's acellular pertussis vaccine was used in Japan starting in 1981. Later versions of the acellular vaccine in other countries consisted of additional defined components of B. pertussis and were often part of the DTaP combination vaccine.
| Biology and health sciences | Infectious disease | null |
170939 | https://en.wikipedia.org/wiki/Step%20function | Step function | In mathematics, a function on the real numbers is called a step function if it can be written as a finite linear combination of indicator functions of intervals. Informally speaking, a step function is a piecewise constant function having only finitely many pieces.
Definition and first consequences
A function is called a step function if it can be written as
, for all real numbers
where , are real numbers, are intervals, and is the indicator function of :
In this definition, the intervals can be assumed to have the following two properties:
The intervals are pairwise disjoint: for
The union of the intervals is the entire real line:
Indeed, if that is not the case to start with, a different set of intervals can be picked for which these assumptions hold. For example, the step function
can be written as
Variations in the definition
Sometimes, the intervals are required to be right-open or allowed to be singleton. The condition that the collection of intervals must be finite is often dropped, especially in school mathematics, though it must still be locally finite, resulting in the definition of piecewise constant functions.
Examples
A constant function is a trivial example of a step function. Then there is only one interval,
The sign function , which is −1 for negative numbers and +1 for positive numbers, and is the simplest non-constant step function.
The Heaviside function , which is 0 for negative numbers and 1 for positive numbers, is equivalent to the sign function, up to a shift and scale of range (). It is the mathematical concept behind some test signals, such as those used to determine the step response of a dynamical system.
The rectangular function, the normalized boxcar function, is used to model a unit pulse.
Non-examples
The integer part function is not a step function according to the definition of this article, since it has an infinite number of intervals. However, some authors also define step functions with an infinite number of intervals.
Properties
The sum and product of two step functions is again a step function. The product of a step function with a number is also a step function. As such, the step functions form an algebra over the real numbers.
A step function takes only a finite number of values. If the intervals for in the above definition of the step function are disjoint and their union is the real line, then for all
The definite integral of a step function is a piecewise linear function.
The Lebesgue integral of a step function is where is the length of the interval , and it is assumed here that all intervals have finite length. In fact, this equality (viewed as a definition) can be the first step in constructing the Lebesgue integral.
A discrete random variable is sometimes defined as a random variable whose cumulative distribution function is piecewise constant. In this case, it is locally a step function (globally, it may have an infinite number of steps). Usually however, any random variable with only countably many possible values is called a discrete random variable, in this case their cumulative distribution function is not necessarily locally a step function, as infinitely many intervals can accumulate in a finite region.
| Mathematics | Specific functions | null |
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